Lecture Notes on Mathematics for Economists 1

by
Takashi Kunimoto

First Version: August 9, 2007

This Version: May 18, 2010

Summer 2010, Department of Economics, McGill University

August 16 - 27 (tentative): Monday - Friday, 10:00am - 1:00pm; at TBA

Instructor: Takashi Kunimoto

Email: takashi.kunimoto@mcgill.ca
Class Web: http://people.mcgill.ca/takashi.kunimoto/?View=Publications
Office: Leacock 438

COURSE DESCRIPTION: This course is designed to provide you with

mathematical tools that are extensively used in graduate economics courses. The topics
which will be covered is

• Sets and Functions;

• Topology in the Euclidean Space;

• Linear Algebra

• Multivariate Calculus;

• Static Optimization;

• (Optional) Correspondences and Fixed Points; and

• (Optional) The First-Order Differential Equations in One Variable.

A good comprehension of the material covered in the notes is essential for successful
graduate studies in economics. Since we are seriously time constrained – which you
might not believe –, it would be very useful for you to carry one of the books provided
below as a reference after you start to study in graduate school in September.

READING:

The main textbook of this course is “Further Mathematics for Economic Analysis.” I
mostly use this book for the course. However, if you don’t find the main textbook helpful
enough, I strongly recommend you that you should buy at least one of the other books
I listed below as well as “Further Mathematics for Economic Analysis.” Of course, you
can buy any math book which you find useful.
1
I am thankful to the students for their comments, questions, and suggestions. Yet, I believe that
there are still many errors in this manuscript. Of course, all remaining ones are my own.

1
 • “Further Mathematics for Economic Analysis,” by Knut Sydsaeter, Peter Ham-
mond, Atle Seierstad, and Atle Strom, Prentice Hall, 2005 (Main Textbook. If
you don’t have any math book or are not confident about your math skill, this
book will help you a lot.)

• “Mathematical Appendix,” in Advanced Microeconomic Theory, Second Edition,

by Geoffrey A. Jehle and Philip J. Reny, (2000), Addison Wesley (Supplementary.
This is the main textbook for Econ 610 but the mathematical appendix of this
book is too concise in many times)

• “Mathematics for Economists,” by Simon and Blume, Norton, (1994). (Supple-

mentary. This book is a popular math book in many Ph.D programs in economics.
There has to be a reason for that, although I don’t know the true one.)

• “Fundamental Methods of Mathematical Economics,” by A. Chiang, McGraw-Hill.

(more elementary and supplementary)

• “Introductory Real Analysis,” by A. N. Kolmogorov and S.V. Fomin, Dover Publi-

cations (very very advanced and supplementary. If you really like math, this is the
book for you.)

OFFICE HOURS: Wednesday and Friday, 2:00pm - 3:00pm

PROBLEM SETS: There will be several problem sets. Problem sets are essential to
help you understand the course and to develop your skill to analyze economic problems.

ASSESSMENT: No grade will be assigned. However, you are expected to do the

problem sets assigned.

2
Contents

1 Introduction 6

2 Preliminaries 9
2.1 Logic . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 9
2.1.1 Necessity and Sufficiency . . . . . . . . . . . . . . . . . . . . . . . 9
2.1.2 Theorems and Proofs . . . . . . . . . . . . . . . . . . . . . . . . . 10
2.2 Set Theory . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 11
2.3 Relations . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 13
2.3.1 Preference Relations . . . . . . . . . . . . . . . . . . . . . . . . . . 13
2.4 Functions . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 14
2.4.1 Least Upper Bound Principle . . . . . . . . . . . . . . . . . . . . . 15

3 Topology in Rn 17
3.1 Sequences on R . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 17
3.1.1 Subsequences . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 18
3.1.2 Cauchy Sequences . . . . . . . . . . . . . . . . . . . . . . . . . . . 19
3.1.3 Upper and Lower Limits . . . . . . . . . . . . . . . . . . . . . . . . 20
3.1.4 Infimum and Supremum of Functions . . . . . . . . . . . . . . . . 21
3.1.5 Indexed Sets . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 21
3.2 Point Set Topology in Rn . . . . . . . . . . . . . . . . . . . . . . . . . . . 21
3.3 Topology and Convergence . . . . . . . . . . . . . . . . . . . . . . . . . . 23
3.4 Properties of Sequences in Rn . . . . . . . . . . . . . . . . . . . . . . . . . 24
3.5 Continuous Functions . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 26

4 Linear Algebra 29
4.1 Basic Concepts in Linear Algebra . . . . . . . . . . . . . . . . . . . . . . . 29
4.2 Determinants and Matrix Inverses . . . . . . . . . . . . . . . . . . . . . . 32
4.2.1 Determinants . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 32
4.2.2 Matrix Inverses . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 32
4.2.3 Cramer’s Rule . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 33
4.3 Vectors . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 34
4.4 Linear Independence . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 35
4.4.1 Linear Dependence and Systems of Linear Equations . . . . . . . . 36
4.5 Eigenvalues . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 37

3
 4.5.1 Motivations . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 37
4.5.2 How to Find Eigenvalues . . . . . . . . . . . . . . . . . . . . . . . 39
4.6 Diagonalization . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 40
4.7 Quadratic Forms . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 42
4.8 Appendix 1: Farkas Lemma . . . . . . . . . . . . . . . . . . . . . . . . . . 44
4.8.1 Preliminaries . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 45
4.8.2 Fundamental Theorem of Linear Algebra . . . . . . . . . . . . . . 46
4.8.3 Linear Inequalities . . . . . . . . . . . . . . . . . . . . . . . . . . . 47
4.8.4 Non-Negative Solutions . . . . . . . . . . . . . . . . . . . . . . . . 47
4.8.5 The General Case . . . . . . . . . . . . . . . . . . . . . . . . . . . 49
4.9 Appendix 2: Linear Spaces . . . . . . . . . . . . . . . . . . . . . . . . . . 49
4.9.1 Number Fields . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 49
4.9.2 Definitions . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 51
4.9.3 Bases, Components, Dimension . . . . . . . . . . . . . . . . . . . . 52
4.9.4 Subspaces . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 53
4.9.5 Morphisms of Linear Spaces . . . . . . . . . . . . . . . . . . . . . . 54

5 Calculus 55
5.1 Functions of a Single Variable . . . . . . . . . . . . . . . . . . . . . . . . . 55
5.2 Real-Valued Functions of Several Variables . . . . . . . . . . . . . . . . . 56
5.3 Gradients . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 56
5.4 The Directional Derivative . . . . . . . . . . . . . . . . . . . . . . . . . . . 57
5.5 Convex Sets . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 58
5.5.1 Upper Contour Sets . . . . . . . . . . . . . . . . . . . . . . . . . . 59
5.6 Concave and Convex Functions . . . . . . . . . . . . . . . . . . . . . . . . 59
5.7 Concavity/Convexity for C 2 Functions . . . . . . . . . . . . . . . . . . . . 60
5.7.1 Jensen’s Inequality . . . . . . . . . . . . . . . . . . . . . . . . . . . 63
5.8 Quasiconcave and Quasiconvex Functions . . . . . . . . . . . . . . . . . . 64
5.9 Total Differentiation . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 68
5.9.1 Linear Approximations and Differentiability . . . . . . . . . . . . . 69
5.10 The Inverse of a Transformation . . . . . . . . . . . . . . . . . . . . . . . 72
5.11 Implicit Function Theorems . . . . . . . . . . . . . . . . . . . . . . . . . . 73

6 Static Optimization 77
6.1 Unconstrained Optimization . . . . . . . . . . . . . . . . . . . . . . . . . . 77
6.1.1 Extreme Points . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 77
6.1.2 Envelope Theorems for Unconstrained Maxima . . . . . . . . . . . 78
6.1.3 Local Extreme Points . . . . . . . . . . . . . . . . . . . . . . . . . 79
6.1.4 Necessary Conditions for Local Extreme Points . . . . . . . . . . . 80
6.2 Constrained Optimization . . . . . . . . . . . . . . . . . . . . . . . . . . . 81
6.2.1 Equality Constraints: The Lagrange Problem . . . . . . . . . . . . 81
6.2.2 Lagrange Multipliers as Shadow Prices . . . . . . . . . . . . . . . . 84
6.2.3 Tangent Hyperplane . . . . . . . . . . . . . . . . . . . . . . . . . . 84
6.2.4 Local First-Order Necessary Conditions . . . . . . . . . . . . . . . 85

4
 6.2.5 Second-Order Necessary and Sufficient Conditions for Local Ex-
treme Points . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 85
6.2.6 Envelope Result for Lagrange Problems . . . . . . . . . . . . . . . 87
6.3 Inequality Constraints: Nonlinear Programming . . . . . . . . . . . . . . . 88
6.4 Properties of the Value Function . . . . . . . . . . . . . . . . . . . . . . . 90
6.5 Constraint Qualifications . . . . . . . . . . . . . . . . . . . . . . . . . . . 91
6.6 Nonnegativity Constraints . . . . . . . . . . . . . . . . . . . . . . . . . . . 94
6.7 Concave Programming Problems . . . . . . . . . . . . . . . . . . . . . . . 95
6.8 Quasiconcave Programming . . . . . . . . . . . . . . . . . . . . . . . . . . 96
6.9 Appendix: Linear Programming . . . . . . . . . . . . . . . . . . . . . . . . 96

7 Differential Equations 97
7.1 Introduction . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 97

8 Fixed Point Theorems 98

8.1 Banach Fixed Point Theorem . . . . . . . . . . . . . . . . . . . . . . . . . 98
8.2 Brouwer Fixed Point Theorem . . . . . . . . . . . . . . . . . . . . . . . . 99

9 Topics on Convex Sets 102

9.1 Separation Theorems . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 102
9.2 Polyhedrons and Polytopes . . . . . . . . . . . . . . . . . . . . . . . . . . 104
9.3 Dimension of a Set . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 105
9.4 Properties of Convex Sets . . . . . . . . . . . . . . . . . . . . . . . . . . . 105

5
Chapter 1

Introduction

I start my lecture with Rakesh Vohra’s message about what economic theory is. He is a
professor at Northwestern University. 1

All of economic theorizing reduces, in the end, to the solution of one of

three problems.

Given a function f and a set S:

1. Find an x such that f (x) is in S. This is the feasibility question.
2. Find an x in S that optimizes f (x). This is the problem of optimality.
3. Find an x in S such that f (x) = x, this is the fixed point problem.
These three problems are, in general, quite difficult. However, if one is
prepared to make assumptions about the nature of the underlying function
(say it is linear, convex or continuous) and the nature of the set S (convex,
compact etc.) it is possible to provide answers and very nice ones at that.

I think this is the biggest picture of economic theory you could have as you go along
this course. Whenever you are at a loss, please come back to this message.

We build our theory on individuals. Assume that all commodities are traded in the
centralized markets. Throughout Econ 610 and 620, we assume that each individual
(consumer and firm) takes prices as given. We call this the price taking behavior assump-
tion. You might ask why individuals are price takers. My answer would be “why not?”
Let us go as far as we can with this behavioral assumption and thereafter try to see the
limitation of the assumption. However, you have to wait for Econ 611 and 621 for how
to relax this assumption. So, stick with this assumption. For each consumer, we want
to know

1. What is the set of “physically feasible bundles? Is there any such a bundle at all
(feasibility)? We call this set the consumption set.
1
See Preface of Advanced Mathematical Economics by Rakesh V. Vohra.

6
 2. What is the set of “financially feasible bundles? Is there any such a bundle at all
(feasibility)? We call this set the budget set.

3. What is the best bundle to the consumer among all feasible bundles (optimality)?
We call this bundle the consumer’s demand.

We can make the exact parallel argument for the firm. What is the set of “technically”
feasible inputs (feasibility)? We call this the production set of the firm. What is the
best combination of inputs to maximize its profit (optimality)? We call this the firm’s
supply. Once we figure out what are feasible and best choices to each consumer and each
firm under any possible circumstance, we want to know if there is any coherent state of
affairs where everybody makes her best choice. In particular, all markets must clear. We
call this coherent state “competitive (Walrasian) equilibrium.” (a fixed point).

If we move from microeconomics to macroeconomics, we must pay special attention

to time. Now, each individual’s budget set does depend upon time. At each point in
time, he can change his asset portfolio so that he smoothes out his consumption plan
and/or production plan over time. If you know exactly when you die, there is no problem.
Because you just leave no money when you die, unless you want to leave some money
to your kids (i.e., altruistic preferences). This is called the finite time horizon problem.
What if you might live longer than you expected with no money left? Then, what do
you do? So, in reality, you don’t know exactly when you die. This situation can be
formulated as the infinite time horizon problem. Do you see why? To deal with the
infinite horizon problem, we use the transversality condition as the terminal condition of
the feasible set. Moreover, his optimization must be taken into account time. You can
also analogously define a sequence of competitive equilibria of the economy.

How can we summarize what we discussed above? Given a (per capita) consump-
tion stream {ct }∞ ∞
t=0 , (per capita) capital accumulation {kt }t=0 , (per capita) GDP stream
∞
{f (kt )}t=0 , capital depreciation rate δ, the population growth rate n, (per capita) con-
sumption growth g, instantaneous utility function of the representative consumer u(·),
the effective discount rate of the representative consumer β > 0, (per capita) wage profile
{wt }∞ ∞
t=0 , and capital interest rate profile {rt }t=0 :

1. Find a {ct }∞
t=0 such that k̇t = f (kt ) − ct − (δ + g + n)kt holds at each t ≥ 1 and
k0 > 0 is exogenously given. This is the feasibility question. Any such {ct } is called
a feasible consumption stream.
 ∞ −βt
2. Find a feasible consumption stream {ct }∞ t=0 that maximizes V0 = 0 e u(ct )dt.
This is the problem of optimality. I assume that V0 < ∞.

3. Find a {rt , wt }∞
t=0 such that V0 (the planner’s optimization) is sustained through
market economies where k̇t = (rt − n − g)kt + wt − ct holds at each t ≥ 1 and
another condition limt→∞ λt e−βt kt = 0 holds. This latter condition is sometimes
called the transversality condition. This is the fixed point problem. This, in fact,
 
can be done by choosing rt = f (kt ) − δt and wt = f (kt ) − f (kt )kt at each t ≥ 1.

7
 With appropriate re-interpretations, the above is exactly what we had in the begin-
ning except the transversality condition, which is a genuine feature of macroeconomics.

8
Chapter 2

Preliminaries

2.1 Logic
Theorems provide a compact and precise format for presenting the assumptions and
important conclusions of sometimes lengthy arguments, and so help identify immediately
the scope and limitations of the result presented. Theorems must be proved and a proof
consists of establishing the validity of the statement in the theorem in a way that is
consistent with the rules of logic.

2.1.1 Necessity and Sufficiency

Consider any two statements, p and q. When we say “p is necessary for q,” we mean
that p must be true for q to be true. For q to be true requires p to be true, so whenever
q is true, we know that p must also be true. So we might have said, instead, that “p is
true if q is true,” or simply that “p is implied by q” (p ⇐ q).

Suppose we know that “p ⇐ q” is a true statement. What if p is not true? Because

p is necessary for q, when p is not true, then q cannot be true, either. But doesn’t this
just say that “q not true” is necessary for “p not true”? Or that “not-q” is implied by
“not-p” (¬q ⇐ ¬p). This latter form of the original statement is called the contrapositive
form.

Let’s consider a simple illustration of these ideas. Let p be the statement, “x is an

integer less than 10.” Let q be the statement that “x is an integer less than 8.” Clearly,
p is necessary for q (q ⇒ p). If we form the contrapositive of these two statements, the
statement ¬p becomes “x is not an integer less than 10,” and “x is not an integer less
than 8.” Then, observe that ¬q ⇐ ¬p. However, ¬p ⇐ ¬q is false. The value of x could
well be 9.

The notion of necessity is distinct from that of sufficiency. When we say “p is

sufficient for q,” we mean that whenever p holds, q must hold. We can say, “p is true
only if q is true,” or that “p implies q” (p ⇒ q). Once again, whenever the statement

9
“p ⇒ q” is true, the contrapositive statement, “¬q ⇒ ¬p” is also true.

Two implications, “p ⇐ q” and “p ⇒ q,” can both be true. When this is so, I say
that “p is necessary and sufficient for q,” or “p is true if and only if q is true,” or “p iff
q.” When “p is necessary and sufficient for q,” we say that the statements p and q are
equivalent and write “p ⇔ q.”

To illustrate briefly, suppose that p and q are the following statements:

• p ≡ “X is yellow,”
• q ≡ “X is a lemon.”
Certainly, if X is a lemon, then X is yellow. Here, p is necessary for q. At the same
time, just because X is yellow does not mean that it must be a lemon. It could be a
banana. So p is not sufficient for q.

2.1.2 Theorems and Proofs

Mathematical theorems usually have the form of an implication or an equivalence, where
one or more statements are alleged to be related in particular ways. Suppose we have
the theorem “p ⇒ q.” Here, p is the assumption and q is the conclusion. To prove a
theorem is to establish the validity of its conclusion given the truth of its assumption,
and several methods can be used to do that.

1. In a constructive proof, we assume that p is true, deduce various consequences of

that, and use them to show that q must also hold. This is also sometimes called a
direct proof, for obvious reasons.
2. In a contrapositive proof, we assume that q does not hold, then show that p cannot
hold. This approach takes advantage of the logical equivalence between the claims
“p ⇒ q” and “¬q ⇒ ¬p” noted earlier, and essentially involves a constructive proof
of the contrapositive to the original statement.
3. In a proof by contradiction, the strategy is to assume that p is true, assume that q
is not true, and attempt to derive a logical contradiction. This approach relies on
the fact that p ⇒ q or ¬q always is true and if “p ⇒ ¬q” is false, then p ⇒ q must
be true.
4. In a proof by mathematical induction, I have a statement H(k) which does depend
upon a natural number k. What I want to show is that a statement H(k) is true
for each k = 1, 2, . . . . First, I show that H(1) is true. This step is usually easy to
establish. Next, I show that H(k) =⇒ H(k + 1), i.e., if H(k) is true, then H(k + 1)
is also true. These two steps allows me to claim that I am done.

If I assert that p is necessary and sufficient for q, or that “p ⇔ q,” we must give a
proof in “both directions.” That is, both “p ⇒ q” and “q ⇒ p” must be established
before a complete proof of the assertion has been achieved.

10
 It is important to keep in mind the old saying that goes, “Proof by example is no
proof.” Suppose the following two statements are given:

• p ≡ “x is a student,”

• q ≡ “x has red hair.”

Assume further that we make the assertion “p ⇒ q.” Then clearly finding one student
with red hair and pointing him out to you is not going to convince you of anything.
Examples are good for illustrating but typically not for proving.

Finally, a sort of converse to the old saying about examples and proofs should be
noted. Whereas citing a hundred examples can never prove that a certain property
always holds, citing one solitary counterexample can disprove that the property always
holds. For instance, to disprove the assertion about the color of students’ hair, you need
simply point out one student with brown hair. A counterexample proves that the claim
cannot always be true because you have found at least one case where it is not.

2.2 Set Theory

A set is any collection of elements. Sets of objects will usually be denoted by capital
letters, A, S, T for example, while their members by lower case, a, s, t for example (English
or Greek). A set S is a subset of another set T if every element of S is also an element
of T . We write S ⊂ T . If S ⊂ T , then x ∈ S ⇒ x ∈ T . The set S is a proper subset of T
if S ⊂ T and S
= T ; sometimes one writes S  T in this case. Two sets are equal sets if
they each contain exactly the same elements. We write S = T whenever x ∈ S ⇒ x ∈ T
and x ∈ T ⇒ x ∈ S. The number of elements in a set S, its cardinality, is denoted
|S|. The upside down “A,” ∀, means “for all,” while the backward “E,” ∃ means “there
exists.”

A set S is empty or is an empty set if it contains no elements at all. It is a subset

of “every” set. For example, if A = {x| x2 = 0, x > 1}, then A is empty. We denote
the empty set by the symbol ∅. The complement of a set S in a universal set U is the
set of all elements in U that are not in S and is denoted S c . For any two sets S and T
in a universal set U , we define the set difference denoted S\T , as all elements in the set
S that are not elements of T . Thus, we can think S c = U \S. The symmetric difference
ST = (S\T ) ∪ (T \S) is the set of all elements that belong to exactly one of the sets S
and T . Note that if S = T , then ST = ∅.

For two sets S and T , we define the union of S and T as the set S∪T ≡ {x| x ∈S or x ∈ T }.
We define the intersection of S and T as the set S ∩ T ≡ {x| x ∈ S and x ∈ T }. Let
. }, we can write
Λ ≡ {1, 2, 3, . . . } be an index set. In stead of writing {S1 , S2 , S3 , . .
{Sλ }λ∈Λ . We would denote the union of all  sets in the collection by λ∈Λ Sλ , and the
intersection of all sets in the collection as λ∈Λ Sλ .

11
 The following are some important identities involving the operations defined above.

• A ∪ B = B ∪ A, (A ∪ B) ∪ C = A ∪ (B ∪ C), A ∪ ∅ = A

• A ∩ B = B ∩ A, (A ∩ B) ∩ C = A ∩ (B ∩ C), A ∩ ∅ = ∅

• A ∪ (B ∩ C) = (A ∪ B) ∩ (A ∪ C), A ∩ (B ∪ C = (A ∩ B) ∪ (A ∩ C) (Distribute laws)

• A\(B ∪ C) = (A\B) ∩ (A\C), A\(B ∩ C) = (A\B) ∪ (A\C) (De Morgan’s laws)

• AB = BA, (AB)C = A(BC), A∅ = A

The collection of all subsets of a set A is also a set, called the power set of A and
denoted by P(A). Thus, B ∈ P(A) ⇐⇒ B ⊂ A.

Example 2.1 Let A = {a, b, c}. Then, P(A) = {∅, {a}, {b}, {c}, {a, b}, {a, c}, {b, c}, {a, b, c}}.
The previous argument reveals its stance that the order of the elements in a set
specification does not matter. In particular, {a, b} = {b, a}. However, on many occasions,
one is interested in distinguishing between the first and the second elements of a pair.
One such example is the coordinates of a point in the x − y-plane. These coordinates
are given as an ordered pair (a, b) of real numbers. The important property of ordered
pairs is that (a, b) = (c, d) if and only if a = c and b = d. The product of two sets S and
T is the set of “ordered pairs” in the form (s, t), where the first element in the pair is a
member of S and the second is a member of T . The product of S and T is denoted

S × T ≡ {(s, t)| s ∈ S, t ∈ T }.

The set of real numbers is denoted by the special symbol R and is defined as

R ≡ {x| − ∞ < x < ∞}.

Any n-tuple, or vector, is just an n dimensional ordered tuple (x1 , . . . , xn ) and can
be thought of as a “point” in n dimensional Euclidean space. This space is defined as
the set product

Rn ≡ R · · × R ≡ {(x1 , . . . , xn ) | xi ∈ R, i = 1, . . . , n}.
 × ·
n times

Often, we want to restrict our attention to a subset of Rn , called the “nonnegative

orthant” and denoted Rn+ , where

Rn+ ≡ {(x1 , . . . , xn ) | xi ≥ 0, i = 1, . . . , n} ⊂ Rn .

Furthermore, we sometimes talk about the strictly “positive orthant” of Rn

Rn++ ≡ {(x1 , . . . , xn ) | xi > 0, i = 1, . . . , n} ⊂ Rn+ .

12
2.3 Relations
Any ordered pair (s, t) associates an element s ∈ S to an element t ∈ T . Any collection
of ordered pairs is said to constitute a binary relation between the sets S and T . Many
familiar binary relations are contained in the product of one set with itself. For example,
let X be the closed unit interval, X = [0, 1]. Then the binary relation ≥ consists of all
ordered pairs of numbers in X where the first one in the pair is greater than or equal to
the second one. When, as here, a binary relation is a subset of the product of one set
X with itself, we say that it is a relation on the set X. A binary relation R on X is
represented by the subset of X × X, i.e., R ⊂ X × X. We can build more structure for
a binary relation on some set by requiring that it possesses certain properties.

Definition 2.1 A relation R in X is reflexive if xRx for all x ∈ X.

For example, ≥ and = on R are reflexive, while > is not.

Definition 2.2 A relation R on X is complete if, for all elements x and y in X, xRy
or yRx.
For example, ≥ on R is complete, while > and = are not. Note that R on X is
reflexive if it is complete.

Definition 2.3 A relation R on X is transitive if, for any three elements x, y, and
z ∈ X, xRy and yRz implies xRz.
For instance, all ≥, =, > on R are transitive.

Definition 2.4 A relation R on X is symmetric if xRy implies yRx and it is anti-

symmetric if xRy and yRx implies x = y.
For example, = is symmetric, while ≥ and > are not. However, ≥ is anti-symmetric,
while > is not as well. A relation R is said to be a partial ordering on X if it is reflexive,
transitive, and anti-symmetric. If a partial ordering is complete, it is called a linear
ordering. For instance, the relation ≥ in R is a linear ordering.

For n ≥ 2, the less-than-or-equal-to relation ≥ on Rn is defined by (x1 , . . . , xn ) ≥

(y1 , . . . , yn ) if and only if xk ≥ yk for k = 1, . . . , n. There is also a strict inequality
relation , which is given by (x1 , . . . , xn )  (y1 , . . . , yn ) if and only if xk > yk for all
k = 1, . . . , n.

2.3.1 Preference Relations

I now talk a little bit about economics. Here I apply the concept of relations to the
consumer choice problem. The number of commodities is finite and equal to n. Each
commodity is measured in some infinitely divisible units. Let x = (x1 , . . . , xn ) ∈ Rn+
be a consumption bundle. Let Rn+ be a consumption set that is the set of bundles the
consumer can conceive. We represent the consumer’s preferences by a binary relation,

13
 
, defined on the consumption set, Rn+ . If x  x , we say that “x is at least as good as

x ,” for this consumer.

Definition 2.5 The binary relation  on X is said to be strict preference relation

  
if, x  x if and only if x  x but x  x.

Definition 2.6 The binary relation ∼ on X is said to be indifference relation if,

  
x ∼ x if and only if x  x and x  x.

Exercise 2.1 Show the following:

1.  on Rn+ is reflexive if it is complete.

2.  on Rn+ is not symmetric.

3. ∼ on Rn+ is symmetric.

2.4 Functions
A function is a relation that associates each element of one set with a single, unique
element of another set. We say that the function f is a mapping, map, or transformation
from one set D to another set R and write f : D → R. We call the set D the domain
and the set R the range of the mapping. If y is the point in the range mapped into by
the point x in the domain, we write y = f (x). In set-theoretic terms, f is a relation from
D to R with the property that for each x ∈ D, there is exactly one y ∈ R such that xf y
(x is related to y via f ).

The image of f is that set of points in the range into which some point in the domain
is mapped, i.e.,

I ≡ {y | y = f (x) for some x ∈ D} ⊂ R.

The inverse image of a set of points S ⊂ I is defined as

f −1 (S) ≡ {x | x ∈ D, f (x) ∈ S} .

The graph of the function f is the set of ordered pairs

G ≡ {(x, y) | x ∈ D, y = f (x)} .

If f (x) = y, one also writes x → y. The squaring function s : R → R, for example,

can then be written as s : x → x2 . Thus, → indicates the effect of the function on an
element of the domain. If f : A → B is a function and S ⊂ A, the restriction of f to S
is the function f |S defined by f |S (x) = f (x) for every x ∈ S. There is nothing in the
definition of a function that prohibits more than one element in the domain from being
mapped into the same element in the range. If, however, every point in the range is
assigned to “at most” a single point in the domain, the function is said to be one-to-one,

14
   
that is, for all x, x ∈ D, whenever f (x) = f (x ), then x = x . If the image is equal to
the range - if for every y ∈ R, there is x ∈ D such that f (x) = y, the function is said
to be onto. If a function is one-to-one and onto (sometimes called bijective), then an
inverse function f −1 : R → D exists that is also one-to-one and onto. The composition
of a function f : A → B and a function g : B → C is the function g ◦ f : A → C given
by (g ◦ f )(a) = g(f (a)) for all a ∈ A.

Exercise 2.2 Show that f (x) = x2 is not a one-to-one mapping.

2.4.1 Least Upper Bound Principle

A set S of real numbers is bounded above if there exists a real number b such that b ≥ x
for all x ∈ S. This number b is called an upper bound for S. A set that is bounded above
has many upper bounds. A least upper bound for the set S is a number b∗ that is an
upper bound for S and is such that b∗ ≤ b for every upper bound b. The existence of a
least upper bound is a basic and non-trivial property of the real number system.

Fact 2.1 (Least Upper Bound Principle) Any nonempty set of real numbers that is
bounded above has a least upper bound.
This principle is rather an axiom of real numbers. A set S can have at most one least
upper bound, because if b∗1 and b∗2 are both least upper bounds for S, then b∗1 ≤ b∗2 and
b∗2 ≤ b∗1 , which thus implies that b∗1 = b∗2 . The least upper bound b∗ of S is often called
the supremum of S. We write b∗ = sup S and b∗ = supx∈S x.

Example 2.2 The set S = (0, 5), consisting of all x such that 0 < x < 5, has many
upper bounds, some of which are 100, 6.73, and 5. Clearly no number smaller than 5
can be an upper bound, so 5 is the least upper bound. Thus, sup S = 5.
A set S is bounded below if there exists a real number a such that x ≥ a for all x ∈ S.
The number a is a lower bound for S. A set S that is bounded below has a greatest
lower bound a∗ , with the property a∗ ≤ x for all x ∈ S, and a∗ ≥ a for all lower bounds
a. The number a∗ is called the infimum of S and we write a∗ = inf S or a∗ = inf x∈S x.
Thus, we summarize

• sup S = the least number greater than or equal to all numbers in S; and

• inf S = the greatest number less than or equal to all numbers in S.

Theorem 2.1 Let S be a set of real numbers and b∗ a real number. Then sup S = b∗ if
and only if the following two conditions are satisfied:

1. x ≤ b∗ for all x ∈ S.

2. For each ε > 0, there exists an x ∈ S such that x > b∗ − ε.

15
 Proof of Theorem 2.1: (=⇒) Since b∗ is an upper bound for S, by definition,
property 1 holds, that is, x ≤ b∗ for all x ∈ S. Suppose, on the other hand, that there
is some ε > 0 such that x ≤ b∗ − ε for all x ∈ S. Define b∗∗ = b∗ − ε. This implies
that b∗∗ is also an upper bound for S and b∗∗ < b∗ . This contradicts our hypothesis that
b∗ is a least upper bound for S. (⇐=) Property 1 says that b∗ is an upper bound for
S. Suppose, on the contrary, that b∗ is not a least upper bound. That is, there is some
other b such that x ≤ b < b∗ for all x ∈ S. Define ε = b∗ − b. Then, we obtain that
x ≤ b∗ − ε for all x ∈ S. This contradicts property 2. 

16
Chapter 3

Topology in Rn

3.1 Sequences on R
A sequence is a function k → x(k) whose domain is the set {1, 2, 3, . . . } of all pos-
itive integers. I denote the set of natural numbers by N = {1, 2, . . . }. The terms
x(1), x(2), . . . , x(k), . . . of the sequence are usually denoted by using subscripts: x1 , x2 , . . . , xk , . . . .
We shall use the notation {xk }∞ k=1 , or simply {xk }, to indicate an arbitrary sequence of
real numbers. A sequence {xk } of real numbers is said to be

1. nondecreasing if xk ≤ xk+1 for k = 1, 2, . . .

2. strictly increasing if xk < xk+1 for k = 1, 2, . . .

3. nonincreasing if xk ≥ xk+1 for k = 1, 2, . . .

4. strictly decreasing if xk > xk+1 for k = 1, 2, . . .

A sequence that is nondecreasing or nonincreasing is called monotone. A sequence

{xk } is said to converge to a number x if xk becomes arbitrarily close to x for all
sufficient large k. We write limk→∞ xk = x or xk → x as k → ∞. The precise definition
of convergence is as follows:

Definition 3.1 The sequence {xk } converge to x if for every ε > 0, there exists a
natural number Nε such that |xk − x| < ε for all k > Nε . The number x is called
the limit of the sequence {xk }. A convergent sequence is one that converges to some
number.
Note that the limit of a convergent sequence is unique. A sequence that does not
converge to any real number is said to diverge. In some cases we use the notation
limk→∞ xk even if the sequence {xk } is divergent. For example, we say that xk → ∞ as
k → ∞. A sequence {xk } is bounded if there exists a number M such that |xk | ≤ M for
all k = 1, 2, . . . . It is easy to see that every convergent sequence is bounded: If xk → x,
by the definition of convergence, only finitely many terms of the sequence can lie outside
the interval I = (x − 1, x + 1). The set I is bounded and the finite set of points from the

17
sequence that are not in I is bounded, so {xk } must be bounded. On the other hand, is
every bounded sequence convergent? No. For example, the sequence {yk } = {(−1)k } is
bounded but not convergent.

Theorem 3.1 Every bounded monotone sequence is convergent.

Proof of Theorem 3.1: Suppose, without loss of generality, that {xk } is nonde-
creasing and bounded. Let b∗ be the least upper bound of the set X = {xk |k = 1, 2, . . . },
and let ε > 0 be an arbitrary number. Theorem 2.1 already showed that there must be a
term xN of the sequence for which xN > b∗ − ε. Because the sequence is nondecreasing,
b∗ − ε < xN ≤ xk for all k > N . But the xk are all less than or equal to b∗ because of
boundedness, so we have b∗ − ε < xk ≤ b∗ . Thus, for any ε > 0, there exists a number
N such that |xk − b∗ | < ε for all k > N . Hence, {xk } converges to b∗ . 

Theorem 3.2 Suppose that the sequences {xk } and {yk } converge to x and y, respec-
tively. Then,

1. limk→∞ (xk ± yk ) = x ± y

2. limk→∞ (xk · yk ) = x · y

3. limk→∞ (xk /yk ) = x/y, assuming that yk

= 0 for all k and y
= 0.

Exercise 3.1 Prove Theorem 3.2.

3.1.1 Subsequences
Let {xk } be a sequence. Consider a strictly increasing sequence of natural numbers

k1 < k2 < k3 < · · ·

and form a new sequence {yj }∞ j=1 , where yj = xkj for j = 1, 2, . . . . The sequence
{yj }j = {xkj }j is called a subsequence of {xk }.

Theorem 3.3 Every subsequence of a convergent sequence is itself convergent, and has
the same limit as the original sequence.
Proof of Theorem 3.3: It is trivial. 

Theorem 3.4 If the sequence {xk } is bounded, then it contains a convergent subse-
quence.
Proof of Theorem 3.4: Since {xk } is bounded, we can assume that there exists
some M ∈ R such that |xk | ≤ M for all k ∈ N. Let yn = sup{xk |k ≥ n} for n ∈ N.
By construction, {yn } is a nonincreasing sequence because the set {xk |k ≥ n} shrinks
as n increases. The sequence {yn } is also bounded because −M ≤ yn ≤ M . Theorem
3.1 already showed that the sequence {yn } is convergent. Let x = limn→∞ yn . By the

18
definition of yn , we can choose a term xkn from the original sequence {xk } (with kn ≥ n)
satisfying |yn − xkn | < 1/n.

|x − xkn | = |x − yn + yn − xkn | ≤ |x − yn | + |yn − xkn | < |x − yn | + 1/n.

This shows that xkn → x as n → ∞. 

3.1.2 Cauchy Sequences

I have defined the concept of convergence of sequences. Then, a natural question arises as
to how we can check if a given sequence is convergent. The concept of Cauchy sequence,
indeed, enables us to do so.

Definition 3.2 A sequence {xk } of real numbers is called a Cauchy sequence if for
every ε > 0, there exists a natural number Nε such that |xm − xn | < ε for all m, n > Nε .
The theorem below is a characterization of convergent sequences.

Theorem 3.5 A sequence is convergent if and only if it is a Cauchy sequence.

Proof of Theorem 3.5: (=⇒) Suppose that {xk } converges to x. Given ε > 0,
we can choose a natural number N such that |xn − x| < ε/2 for all n > N . Then, for
m, n > N ,

|xm − xn | = |xm − x + x − xn | ≤ |xm − x| + |x − xm | < ε/2 + ε/2 = ε.

Therefore, {xk } is a Cauchy sequence. (⇐=) Suppose that {xk } is a Cauchy sequence.
First, we shall show that the sequence is bounded. By the Cauchy property, there is a
number M such that |xk −xM | < 1 for k > M . Moreover, the finite set {x1 , x2 , . . . , xM −1 }
is clearly bounded. Hence, {xk } is bounded. Theorem 3.4 showed that the bounded
sequence {xk } has a convergent subsequence {xkj }. Let x = limj xkj . Because {xk } is a
Cauchy sequence, for every ε > 0, there is a natural number N such that |xm − xn | < ε/2
for all m, n > N . If we take J is sufficiently large, we have |xkj − x| < ε/2 for all j > J.
Then for k > N and j > max{N, J},

|xk − x| = |xk − xkj + xkj − x| ≤ |xk − xkj | + |xkj − x| < ε/2 + ε/2 = ε

Hence xk → x as k → ∞. 

Exercise 3.2 Consider the sequence {xk } with the generic term

1 k
1 1 1
xk = 2 + 2 + · · · + 2 =
1 2 k i2
i=1

19
Prove that this sequence is a Cauchy sequence. Hint:
1 1 1
2
+ 2
+ ··· +
(n + 1) (n + 2) (n + k)2
1 1 1
< + + ··· +
n(n + 1) (n + 1)(n + 2) (n + k − 1)(n + k)




1 1 1 1 1 1
= − + − + ··· + −
n n+1 n+1 n+2 n+k−1 n+k

Exercise 3.3 Prove that a sequence can have at most one limit. Use proof by contra-
diction. Namely, you first suppose, by way of contradiction, that there are two limit
points.

3.1.3 Upper and Lower Limits

Let {xk } be a sequence that is bounded above, and define yn = sup{xk |k ≥ n} for
n = 1, 2, . . . . Each yn is a finite number and {yn } is a nonincreasing sequence. Then
either limn→∞ yn exists or is −∞. We call this limit the upper limit (or lim sup) of the
sequence {xk }, and we introduce the following notation:

lim sup xk = lim (sup{xk |k ≥ n})

k→∞ n→∞

If {xk } is not bounded above, we write lim supk→∞ xk = ∞. Similarly, if {xk } is

bounded below, its lower limit (or lim inf, is defined as

lim inf xk = lim (inf{xk |k ≥ n)

k→∞ n→∞

If {xk } is not bounded below, we write lim inf k→∞ xk = −∞.

Theorem 3.6 If the sequence {xk } is convergent, then

lim sup xk = lim inf xk = lim xk .

k→∞ k→∞ k→∞

On the other hand, if limk→∞ sup xk = limk→∞ inf xk , then {xk } is convergent.
I omit the proof of Theorem 3.6.

Exercise 3.4 Determine the lim sup and lim inf of the following sequences.

1. {xk } = {(−1)k }

2. {xk } = (−1)k (2 + 1/k) + 1

20
3.1.4 Infimum and Supremum of Functions
Suppose that f (x) is defined for all x ∈ B, where B ⊂ Rn . We define the infimum and
supremum of the function f over B by

inf f (x) = inf{f (x)|x ∈ B}, sup f (x) = sup{f (x)|x ∈ B}.
x∈B x∈B

If a function f is defined over a set B, if inf x∈B f (x) = y, and if there exists a c ∈ B
such that f (c) = y, then we say that the infimum is attained (at the point c) in B. In
this case the infimum y is called the minimum of f over B, and we often write “min”
instead of “inf.” In the same way we write “max” instead of “sup” when the supremum
of f over B is attained in B, and so becomes the maximum.

3.1.5 Indexed Sets

Suppose that, for each λ ∈ Λ, we specify an object aλ . Then, these objects form an
indexed set {aλ }λ∈Λ with Λ as its index set. In formal terms, an indexed set is a function
whose domain is the index set. For example, a sequence is an indexed set {ak }k∈Æ with
the set N of natural numbers as its index set. In stead of {ak }k∈Æ one often writes
{ak }∞
k=1 .

A set whose elements are sets is often called a family of sets, and so an indexed set
of sets is also called an indexed family of sets. Consider a nonempty indexed family
{Aλ }λ∈Λ of sets. The union and the intersection of this family are the sets

• λ∈Λ Aλ = the set consisting of all x that belong to Aλ for at least one λ ∈ Λ

• λ∈Λ Aλ = the set consisting of all x that belong to Aλ for all λ ∈ Λ.

The distributive laws can be generalized to

   
   
A∪ Bλ = (A ∪ Bλ ) , A ∩ Bλ = (A ∩ Bλ )
λ∈Λ λ∈Λ λ∈Λ λ∈Λ

and De Morgan’s laws to

   
   
A\ Bλ = (A\Bλ ) , A\ Bλ = (A\Bλ )
λ∈Λ λ∈Λ λ∈Λ λ∈Λ

The union
∞ and the intersection
∞ of a sequence {An }n∈Æ = {An }∞
n=1 of sets is often
written as n=1 An and n=1 An .

3.2 Point Set Topology in Rn

Consider the n-dimensional Euclidean space Rn , whose elements, or points, are n-
vectors x = (x1 , . . . , xn ). The Euclidean distance d(x, y) between any two points x =

21
(x1 , . . . , xn ) and y = (y1 , . . . , yn ) in Rn is the norm x − y of the vector difference
between x and y. Thus,

d(x, y) = x − y = (x1 − y1 )2 + · · · + (xn − yn )2

If x, y, and z are points in Rn , then

d(x, z) ≤ d(x, y) + d(y, z) (triangle inequality)

If x0 is a point in Rn and r is a positive real number, then the set of all points x ∈ Rn
whose distance from x0 is less than r, is called the open ball around x0 with radius r.
This open ball is denoted by Br (x0 ). Thus,

Br (x0 ) = {x ∈ Rn |d(x0 , x) < r}

Definition 3.3 A set S ⊂ Rn is open if, for all x0 ∈ S, there exists some ε > 0 such
that Bε (x0 ) ⊂ S.
On the real line R, the simplest type of open set is an open interval. Let S be any
subset of Rn . A point x0 ∈ S is called an interior point of S if there is some ε > 0 such
that Bε (x0 ) ⊂ S. The set of all interior points of S is called the interior of S, and is
denoted int(S). A set S is said to be a neighborhood of x0 if x0 is an interior point of S,
that is, if S contains some open ball Bε (x0 ) (i.e., Bε (x0 ) ⊂ S) for some ε > 0.

Theorem 3.7 1. The entire space Rn and the empty set ∅ are both open.

2. arbitrary unions of open sets are open: Let Λ be an arbitrary index set. If Aλ is
open for each λ ∈ Λ, then, λ∈Λ Aλ is also open.

3. The intersection of finitelymany open sets is open: Let Λ be a finite set. If Aλ is

open for each λ ∈ Λ, then λ∈Λ Aλ is open.
Proof of Theorem 3.7: (1) It is clear that B1 (x) ⊂ Rn for all x ∈ Rn , so Rn is
open. The empty set ∅ is open because the set has no element, so every member is an
interior point.

(2) Let {Uλ }λ∈Λ be an arbitrary family of open sets in Rn , and let U ∗ = λ∈Λ Uλ be
the union of the whole family. For each x ∈ U ∗ , there is at least one λ ∈ Λ such that
x ∈ Uλ . Since Uλ is open by our hypothesis, there exists ε > 0 such that Bε (x) ⊂ Uλ ⊂
U ∗ . Hence, U ∗ is open.
K
(3) Let {Ui }K n
k=1 be finite collection of open sets in R , and let U∗ = k=1 Uk be the
intersection of all these sets. Let x be any point in U∗ . Since Uk is open for each k,
there is εk > 0 such that Bεk (x) ⊂ Uk for each k. Let ε∗ = min{ε1 , . . . , εK }. This is
well defined because of finiteness. Then, Bε∗ (x) ⊂ Uk for each k, which implies that
Bε∗ (x) ⊂ U∗ . Hence, U∗ is open. 

22
Exercise 3.5 There are two questions. First, draw the graph of S = {(x, y) ∈ R2 |2x −
y < 2 and x − 3y < 5}. Second, prove that S is open in R2 .

Definition 3.4 A set S is closed if its complement, Rn \S is open.

A point x0 ∈ Rn is said to be a boundary point of the set S ⊂ Rn if Bε (x0 ) ∩ S c
= ∅
and Bε (x0 ) ∩ S
= ∅ for every ε > 0. Here, S c = R\S. In general, a set may include
none, some, or all of its boundary points. An open set, for instance, contains none of its
boundary points.

Each point in a set is either an interior point or a boundary point of the set. The set
of all boundary points of a set S is said to be the boundary of S and is denoted ∂S or
bd(S). Note that, given any set S ⊂ Rn , there is a corresponding partition of Rn into
three mutually disjoint sets (some of which may be empty), namely;

1. the interior of S, which consists of all points x ∈ Rn such that N ⊂ S for some
neighborhood N of x;

2. the exterior of S, which consists of all points x ∈ Rn for which there exists some
neighborhood N of x such that N ⊂ Rn \S;

3. the boundary of S, which consists of all points x ∈ Rn with the property that every
neighborhood N of x intersects both S and its complement Rn \S.

A set S ∈ Rn is said to be closed if it contains all its boundary points. The union of
S and its boundary (S ∪ ∂S) is called the closure of S, denoted by S̄. A point x belongs
to S̄ if and only if Bε (x) ∩ S
= ∅ for any ε > 0. The closure S̄ of any set S is indeed
closed. In fact, S̄ is the smallest closed set containing S.

Theorem 3.8 1. The whole space Rn and the empty set ∅ are both closed.

2. Arbitrary intersections of closed sets are closed.

3. The union of finitely many closed sets is closed.

Exercise 3.6 Prove Theorem 3.8. Use the fact that the complement of open set is closed
and Theorem 3.7.
In topology, any set containing some of its boundary points but not all of them, is
neither open nor closed. The half-open intervals [a, b) and (a, b], for examples, are neither
open nor closed. Hence, openness and closedness are not mutually exclusive.

3.3 Topology and Convergence

I want to generalize the argument in Section 3.1 into Rn . The basic idea to do so is to
apply the previous argument coordinate-wise. A sequence {xk }∞ n
k=1 in R is a function
n
that for each natural number k yields a corresponding point xk in R .

23
Definition 3.5 A sequence {xk } in Rn converges to a point x ∈ Rn if for each ε > 0,
there exists a natural number N such that xk ∈ Bε (x) for all k ≥ N , or equivalently, if
d(xk , x) → 0 as k → ∞.

Theorem 3.9 Let {xk } be a sequence in Rn . Then, {xk } converges to the vector x ∈ Rn
(j)
if and only if for each j = 1, . . . , n, the real number sequence {xk }∞k=1 , consisting of jth
(j)
component of each vector xk , converges to x ∈ R, the jth component of x.
Proof of Theorem 3.9: (=⇒) For every k and every j, one has d(xk , x) = xk −x| ≥
(j) (j)
|xk − x(j) |. It follows
that if xk → x, then xk → x(j) for each j. (⇐=) Suppose that
(j)
xk → x(j) as k → ∞ for j = 1, . . . , n. Then, given any ε > 0, for each i = 1, . . . , n,
(j) √
there exists a number Nj such that |xk − x(j) | < ε/ n for all k > Nj . It follows that
 
(1) (n)
d(xk , x) = |xk − x(1) |2 + · · · + |xk − x(n) |2 < ε2 /n + · · · ε2 /n = ε,

for all k > max{N1 , . . . , Nn }. This is well defined because of the finiteness of n. There-
fore, xk → x as k → ∞ 

Definition 3.6 A sequence {xk } in Rn is said to be Cauchy sequences if for every

ε > 0, there exists a number N such that d(xm , xn ) < ε for all m, n > N .

Theorem 3.10 A sequence {xk } in Rn is convergent if and only if it is a Cauchy se-

quence.

Exercise 3.7 Prove Theorem 3.10. Apply the same argument in Theorem 3.5 to each
coordinate.

3.4 Properties of Sequences in Rn

Theorem 3.11 1. For any set S ⊂ Rn , a point x ∈ Rn belongs to S̄ if and only if
there exists {xk } ∈ S such that xk → x as k → ∞.

2. A set S ⊂ Rn is closed if and only if every convergent sequence of points in S has

its limit in S.
Proof of Theorem 3.11: (=⇒ of Property 1) Let x ∈ S̄. Regardless of whether
x ∈ intS or ∂S, for each k ∈ N, we can construct xk such that xk ∈ B1/k (x) ∩ S (In
particular, take xk = x for each k). Then xk → x as k → ∞. (⇐= of Property 1) Suppose
that {xk } is a convergent sequence for which xk ∈ S for each k and x = limk→∞ xk . We
claim that x ∈ S̄. For any ε > 0, there is a number N ∈ N such that xk ∈ Bε (x) for
all k > N . Since xk ∈ S for each k, it follows that Bε (x) ∩ S
= ∅. Suppose, on the
other hand, x ∈ / S̄. Since S̄ is closed, there is some ε > 0 such that Bε (x) ∩ S = ∅. This
contradicts the conclusion I just drew that Bε (x) ∩ S
= ∅ for any ε > 0. Hence, x ∈ S̄.

24
 (=⇒ of Property 2) Assume that S is closed and let {xk } be a convergent sequence
such that xk ∈ S for each k. Note that x ∈ S̄ by property 1. Since S = S̄ if S is closed,
it follows that x ∈ S. (⇐= of Property 2) By property 1, for any point x ∈ S̄, there is
some sequence {xk } for which xk ∈ S for each k and limk→∞ xk = x. By our hypothesis,
x ∈ S. This shows that x ∈ S̄ implies x ∈ S, i.e., S̄ ⊂ S. By definition, S ⊂ S̄ for any
S. Hence S = S̄, that is, S is closed. 

Definition 3.7 A set S in Rn is bounded if there exists a number M ∈ R such that

x ≤ M  for all x ∈ S. A set that is not bounded is called unbounded. Here x =
d(x, 0) = x21 + · · · + x2n , called the Euclidean norm.
Similarly, a sequence {xk } in Rn is bounded if the set {xk |k = 1, 2, . . . } is bounded.

Lemma 3.1 Any convergent sequence {xk } in Rn is bounded.

Proof of Lemma 3.1: If xk → x, then only finitely many terms of sequence can lie
outside the ball B1 (x). The ball B1 (x) is bounded and any finite set of points is bounded,
so {xk } must be bounded. 

On the other hand, a bounded sequence {xk } in Rn is not necessarily convergent.

This is the same as the sequences in R. The theorem below gives us a characterization
of boundedness of the set in terms of sequences.

Theorem 3.12 A subset S of Rn is bounded if and only if every sequence of points in

S has a convergent subsequence.
I omit the proof of this theorem. Although it is not a difficult proof, it is tedious to
prove that. The next concept of compact sets is used extensively in both mathematics
and economics.

Definition 3.8 A set S in Rn is compact if it is closed and bounded.

Theorem below is a characterization of compact sets in terms of sequences.

Theorem 3.13 (Bolzano-Weierstrass) A subset S of Rn is compact if and only if

every sequence of points in S has a subsequence that converges to a point in S.
Proof of Bolzano-Weierstrass’s theorem: (=⇒) Suppose that S is compact and
let {xk } be a sequence such that xk ∈ S for each k. Since S is also bounded, there is
a convergent subsequence {yn } = {xkn }. Furthermore, limn→∞ yn = y ∈ S because S
is closed. (⇐=) Let {xk } be any sequence for which xk ∈ S for each k. Suppose that
there exists a subsequence {yn } = {xkn } of the sequence {xk } with limn→∞ yn = y ∈ S.
This with the previous theorem (Theorem 3.12) already showed that S is bounded. Let
{xk } be any sequence for which xk ∈ S for each k and x = limk→∞ xk . Since S̄ is closed
by definition, it follows that x ∈ S̄. By assumption, {xk } has a subsequence {xkj } that

converges to limj→∞ xkj = x ∈ S. But {xkj } also converges to x. Hence, the limit

points must be the same, that is, x = x ∈ S. 

25
Exercise 3.8 Let the number of commodities of the competitive market be n. Let pi > 0
be a price for commodity i for each i = 1, . . . , n. Let y > 0 be the consumer’s income.
Define the consumer’s budget set B(p, y) as
  n


B(p, y) ≡ x = (x1 , . . . , xn ) ∈ Rn+  pi xi ≤ y .
i=1

Show that B(p, y) is nonempty and compact.

3.5 Continuous Functions

Consider first a real-valued function z = f (x) = f (x1 , . . . , xn ) of n variables. Roughly
speaking, f is continuous if small changes in the independent variables cause only small
changes in the function value.

Definition 3.9 A function f : S → R with domain S ⊂ Rn is continuous at a point

x0 in S if for every ε > 0 there exists a δ > 0 such that

|f (x) − f (x0 )| < ε for all x ∈ S with x − x0  < δ.

If f is continuous at every point in a set S, we simply say that f is continuous on S.

√
Exercise 3.9 Let f (x) = x be a function f : R+ → R. Prove that f is continuous.

Exercise 3.10 Let f : R+ → R be given below.


1 if x ≥ 1
f (x) =
0 if x < 1

Show that f is not a continuous function.

Consider next the general case of vector-valued functions.

Definition 3.10 A function f = (f1 , . . . , fm ) from a subset S of Rn to Rm is said to be

continuous at x0 in S if for every ε > 0, there exists a δ > 0 such that d(f (x), f (x0 )) < ε
for all x ∈ S with d(x, x0 ) < δ, or equivalently, such that f (Bδ (x0 ) ∩ S) ⊂ Bε (f (x0 )).
The next theorem shows that the continuity of vector-valued functions can reduce to
the continuity of each component (coordinate) functions, and vice versa.

Theorem 3.14 A function f = (f1 , . . . , fm ) from S ⊂ Rn to Rm is continuous at a

point x0 in S if and only if each component function fj : S → R, j = 1, . . . , m, is
continuous at x0 .

26
 Proof of Theorem 3.14: (=⇒) Suppose f is continuous at x0 . Then, for every
ε > 0, there exists a δ > 0 such that

|fj (x) − fj (x0 )| ≤ d(f (x), f (x0 )) < ε

for every x ∈ S with d(x, x0 ) < δ. Hence, fj is continuous at x0 for j = 1, . . . , m. (⇐=)

Suppose that each component fj is continuous at x0 . Then, for every ε > 0 and every
√
j = 1, . . . , m, there exists δj > 0 such that |fj (x)− fj (x0 )| < ε/ m for every point x ∈ S
with d(x, x0 ) < δj . Let δ = min{δ1 , . . . , δm }. Then x ∈ Bδ (x0 ) ∩ S implies that

 ε2 ε2
d(f (x), f (x0 )) = |f1 (x) − f1 (x0 )|2 + · · · + |fm (x) − fm (x0 )|2 < + ··· + = ε.
m m
This proves that f is continuous at x0 . 

Here, I want to characterize the continuity of the functions in terms of sequences.

Theorem 3.15 A function f from S ⊂ Rn into Rm is continuous at a point x0 in S if

and only if f (xk ) → f (x0 ) for every sequence {xk } of points in S that converges to x0 .
Proof of Theorem 3.15: (=⇒) Suppose that f is continuous at x0 and let {xk } be
a sequence for which xk ∈ S and limk→∞ xk = x0 . Let ε > 0 be given. Since xk → x0 ,
for any δ > 0, there exists a number N ∈ N such that d(xk , x0 ) < δ for all k > N .
Therefore, because of the continuity of f , there exists δ > 0 such that d(f (x), f (x0 )) < ε
whenever x ∈ Bδ (x0 ) ∩ S. But then xk ∈ Bδ (x0 ) ∩ S and so d(f (xk ), f (x0 )) < ε for
all k > N . This implies that f (xk ) → f (x0 ). (⇐=) Let {xk } be a sequence for which
xk ∈ S for each k and x0 = limk→∞ xk . Since xk → x0 , for any δ > 0, there is a number
Nδ such that d(xk , x0 ) < δ for all k > Nδ . Similarly, since f (xk ) → f (x0 ), for any
ε > 0, there exists a number Nε such that d(f (xk ), f (x0 )) < ε for all k > Nε . Define
N ∗ = max{Nδ , Nε }. Then, by choosing k > N ∗ , there exists a δ > 0 with d(xk , x0 ) < δ
such that d(f (xk ), f (x0 )) < ε. In fact, this holds for any ε > 0. Hence, we prove that f
is continuous at x0 . 

The theorem below shows that continuous mappings preserve the compactness of the
set.

Theorem 3.16 Let S ⊂ Rn and let f : S → Rm be continuous. Then f (K) = {f (x)|x ∈

K} is compact for every compact subset K of S.
Proof of Theorem 3.16: Let {yk } be any sequence in f (K). By definition, for each
k, there is a point xk ∈ K such that yk = f (xk ). Because K is compact, by Bolzano-
Weierstrass’s theorem, the sequence {xk } has a subsequence {xkj } with the property that
xkj ∈ K for each j and limj→∞ xkj = x0 ∈ K. Because f is continuous, by the previous
theorem (Theorem 3.15), f (xkj ) → f (x0 ) as j → ∞, where f (x0 ) ∈ f (K) because
x0 ∈ K. But then {ykj } is a subsequence of {yk } that converges to a point f (x0 ) ∈ f (K).
So, we have proved that any sequence in f (K) has a subsequence converging to a point
of f (K). 

27
 Suppose that f is a continuous function from Rn to Rm . If V is an open set in Rn ,
the image f (V ) = {f (x)|x ∈ V } of V need not be open in Rm . Nor need f (C) be closed
if C is closed. Nevertheless, the inverse image f −1 (U ) = {x|f (x) ∈ U } of an open set U
under continuous function f is always open. Similarly, the inverse image of any closed
set must be closed.

Theorem 3.17 Let f be any function from Rn to Rm . Then f is continuous if and only
if either of the following equivalent conditions is satisfied.

1. f −1 (U ) is open for each open set U in Rm .

2. f −1 (F ) is closed for each closed set F in Rm .

I omit the proof of Theorem 3.17. Because it is conceptually involved. So, just accept
the result.

Theorem 3.18 Let S be a compact set in R and let x∗ be the greatest lower bound of S
and x∗ be the lowest upper bound of S. Then, x∗ ∈ S and x∗ ∈ S.
Proof of Theorem 3.18: Let S ⊂ R be closed and bounded and let x∗ be the
lowest upper bound of S. Then, by definition of any lower bound, we have x∗ ≥ x for
all x ∈ S. If x∗ = x for some x ∈ S, we are done. Suppose, therefore, that x∗ is strictly
greater than every point in S. If x∗ > x for all x ∈ S, then x∗ ∈ / S, so x∗ ∈ R\S.
Since S is closed, R\S is open. Then, by the definition of open sets, there exists some
ε > 0 such that Bε (x∗ ) = (x∗ − ε, x∗ + ε) ⊂ R\S. Since x∗ > x for all x ∈ S and
Bε (x∗ ) ⊂ R\S, we claim that for any x̃ ∈ Bε (x∗ ), we must have x̃ > x for all x ∈ S. In
particular, x∗ − ε/2 ∈ Bε (x∗ ) and x∗ − ε/2 > x for all x ∈ S. But then this contradicts
our hypothesis that x∗ is the lowest upper bound of S. Thus, we must conclude that
x∗ ∈ S. The same argument should be constructed for the greatest lower bound of S,
i.e., x∗ . 

Theorem 3.19 (Weierstrass’s Theorem) Let f : S → R be a continuous real-valued

mapping where S is a nonempty compact subset of Rn . Then there exists a vector x∗ , x∗ ∈
S such that for all x ∈ S,

f (x∗ ) ≤ f (x) ≤ f (x∗ ).

Proof of Weierstrass’s Theorem: It follows from Theorem 3.16 and 3.18. 

28
Chapter 4

Linear Algebra

4.1 Basic Concepts in Linear Algebra

Let f : Rn → Rm be a mapping (transformation). A mapping f is said to be linear
if for any x, y ∈ Rn and any α ∈ R, the following two conditions are satisfied: (1)
f (x + y) = f (x) + f (y) and (2) f (αx) = αf (x). For any linear mapping f : Rn → Rm ,
there exists a unique m × n matrix A such that f (x) = Ax for all x ∈ Rn . 1 An m × n
matrix is a rectangular array with m rows and n columns:
⎛ ⎞
a11 a12 · · · a1n
⎜ a21 a22 · · · a2n ⎟
⎜ ⎟
A = (aij )m×n = ⎜ . .. . .. ⎟
⎝ . . . . . . ⎠
am1 am2 · · · amn

Here aij denotes the elements in the ith row and the jth column. With this notation,
we can express f (x) = Ax as below:
⎛ (1) ⎞
f (x) ⎛ ⎞⎛ ⎞
⎜ .. ⎟ a11 a12 ··· a1n x1
⎜ . ⎟ ⎜ ⎟⎜ ⎟
⎜ ⎟ ⎜ a21 a22 ··· a2n ⎟⎜ x2 ⎟
f (x) = ⎜ (j) ⎟
⎜ f (x) ⎟ = ⎜ . .. .. .. ⎟⎜ .. ⎟.
⎜ .. ⎟ ⎝ .. . . . ⎠⎝ . ⎠
⎝ . ⎠
am1 am2 · · · amn xn
(m)
f (x)

Exercise 4.1 Show that any linear mapping f : Rn → Rm is continuous.

If A = (aij )m×n , B = (bij )m×n , and α is a scalar, we define

• A + B = (aij + bij )m×n ,

• αA = (αaij )m×n ,
1
This is a non-trivial statement. But I take this one-to-one correspondence between linear mapping
and matrix representation as a fact with no proof provided.

29
 • A − B = A + (−1)B = (aij − bij )m×n .
Let f : Rn → Rm and g : Rn → Rp be linear mappings. Then, we can set m × n
matrix A = (aij )m×n associated with f and p × m matrix B = (bij )p×m associated with
g. Consider the composite mapping g ◦ f (x) = g(f (x)). What I would like to have is the
requirement on the product of matrices that g ◦ f ≡ BA. Then the product C = BA
is defined as the p × n matrix C = (cij )p×n , whose element in the ith row and the jth
column is the inner product of the ith row of A and the jth column of B. That is,
n

cij = air brj = ai1 b1j + ai2 b2j + · · · + aik bkj + · · · ain bnj
r=1
  
n terms

It is important to note that the product BA is well defined only if the number of
columns in B is equal to the number of rows in A.

If A, B, and C are matrices whose dimensions are such that the given operations are
well defined, then the basic properties of matrix of multiplication are:
• (AB)C = A(BC) (associative law)
• A(B + C) = AB + AC (left distributive law)
• (A + B)C = AC + BC (right distributive law)

Exercise 4.2 Show the above three properties when we consider 2 × 2 matrices.
However, matrix multiplication is not commutative. In fact,
• AB
= BA, except in special cases
• AB = 0 does not imply that A or B is 0
• AB = AC and A
= 0 do not imply that B = C

Exercise 4.3 Confirm the above three points by example.

By using matrix multiplication, one can write a general system of linear equations in
a very concise way. Specifically, the system
a11 x1 + a12 x2 + · · · + a1n xn = b1
a21 x1 + a22 x2 + · · · + a2n xn = b2
··························· ··· ···
am1 x1 + am2 x2 + · · · + amn xn = bm
can be written as Ax = b if we define
⎛ ⎞ ⎛ ⎞ ⎛ ⎞
a11 a12 ··· a1n x1 b1
⎜ a21 a22 ··· a2n ⎟ ⎜ x2 ⎟ ⎜ b2 ⎟
⎜ ⎟ ⎜ ⎟ ⎜ ⎟
A=⎜ . .. .. .. ⎟, x = ⎜ .. ⎟, b = ⎜ .. ⎟
⎝ .. . . . ⎠ ⎝ . ⎠ ⎝ . ⎠
am1 am2 ··· amn xn bn

30
 A matrix is square if it has an equal number of rows and columns. If A is a square
matrix and n is a positive integer, we define the nth power of A in the obvious way:

An = AA
 · · · A
n factors

For diagonal matrices it is particularly easy to compute powers:

⎛ ⎞ ⎛ n ⎞
d1 0 ··· 0 d1 0 · · · 0
⎜ 0 d2 · · · 0 ⎟ ⎜ 0 dn · · · 0 ⎟
⎜ ⎟ n ⎜ 2 ⎟
D=⎜ . .. . . .. ⎟ =⇒ D = ⎜ .. .. . . .. ⎟
⎝ .. . . . ⎠ ⎝ . . . . ⎠
0 0 ··· dm 0 0 ··· dnm

The identity matrix of order n, denoted by In , is the n × n matrix having ones along
the main diagonal and zeros elsewhere:
⎛ ⎞
1 0 ··· 0
⎜ 0 1 ··· 0 ⎟
⎜ ⎟
In = ⎜ . . .
. . ... ⎟
(identity matrix)
⎝ .. .. ⎠
0 0 ··· 1

If A is any m × n matrix, then AIn = A = Im A. In particular,

AIn = In A = A for every n × n matrix A

If A = (aij )m×n is any matrix, the transpose of A is defined as AT = (aji )n×m .

The subscripts i and j are interchanged because every row of A becomes a column of
AT , and every column of A becomes a row of AT . The following rules apply to matrix
transposition:

1. (AT )T = A

2. (A + B)T = AT + B T

3. (αA)T = αAT

4. (AB)T = B T AT

Exercise 4.4 Prove the above four properties when we consider 2 × 2 matrices.
A square matrix is said to be symmetric if A = AT .

31
4.2 Determinants and Matrix Inverses
4.2.1 Determinants
Recall that the determinants |A| of 2 × 2 and 3 × 3 matrices are defined by
 
 a a12 
|A| =  11 = a11 a22 − a12 a21
a21 a22 
 
 a11 a12 a13 

|A| =  a21 a22 a23  = a11 a22 a33 + a12 a23 a31 + a13 a21 a32 − a11 a23 a32 − a12 a21 a33 − a13 a22 a31
 a31 a32 a33 
For a general n × n matrix A = {aij }, the determinant |A| can be defined recursively.
In fact,
|A| = ai1 Ai1 + ai2 Ai2 + · · · aij Aij + · · · + ain Ain
where the cofactors Aij are determinants of (n − 1) × (n − 1) matrices given by
 
 a11 · · · a1,j−1 a1j a1,j+1 · · · a1n 

 a21 · · · a2,j−1 a2j a2,j+1 · · · a2n 

 .. .. .. .. .. 

i+j  . . . . . 
Aij = (−1) 
 a i1 · · · ai,j−1 a ij ai,j+1 · · · ain 
 . .. .. .. .. 
 .. . . . . 

 a · · · ann 
n1 · · · an,j−1 anj an,j+1
Here row i and column j are to be deleted from the matrix A to produce Aij .

Proposition 4.1 Let A and B be n × n matrices. Then, |AB| = |A||B|

Exercise 4.5 Prove Proposition 1 when n = 2.

4.2.2 Matrix Inverses

The inverse A−1 of an n × n matrix A has the following properties:
• B = A−1 ⇐⇒ AB = In ⇐⇒ BA = In
• A−1 exists ⇐⇒ |A|
= 0
If A = (aij )n×n and |A|
= 0, the unique inverse of A is given by
⎛ ⎞
A11 A21 · · · An1
⎜ ⎟
1 ⎜ A12 A22 · · · An2 ⎟
A−1 = adj(A), where adj(A) = ⎜ . .. .. .. ⎟
|A| ⎝ .. . . . ⎠
A1n A2n · · · Ann
with Aij , the cofactor of the element aij . Note carefully the order of the indices in the
adjoint matrix, adj(A) with the column number preceding the row number. The matrix
(Aij )n×n is called the cofactor matrix, whose transpose is the adjoint matrix.

32
Exercise 4.6 Define A as
⎛ ⎞
a11 a12 a13
A = ⎝ a21 a22 a23 ⎠
a31 a32 a33

Then, derive A−1 . Assume that |A|

= 0.

Lemma 4.1 The following rules for inverses can be established.

• (A−1 )−1 = A,
• (AB)−1 = B −1 A−1 ,
• (AT )−1 = (A−1 )T ,
• (αA)−1 = α−1 A−1 , where α ∈ R.

Exercise 4.7 Prove Lemma 4.1 when n = 2.

Proposition 4.2 Let A be a n × n matrix. Then, |A−1 | = 1/|A|.

Exercise 4.8 Prove Proposition 4.2 when n = 2.

4.2.3 Cramer’s Rule

A linear system of n equations and n unknowns,

a11 x1 + a12 x2 + · · · + a1n xn = b1

a21 x1 + a22 x2 + · · · + a2n xn = b2 (∗)
··························· ··· ···
an1 x1 + an2 x2 + · · · + ann xn = bn

has a unique solution if and only if |A|

= 0. The solution is then
|Aj |
xj = , j = 1, . . . , n
|A|
where the determinant
 
 a11 · · · a1,j−1 b1 a1,j+1 · · · a1n 
 
 
 a21 · · · a2,j−1 b2 a2,j+1 · · · a2n 
|Aj | =  . .. .. .. .. .. ..


 .. . . . . . . 
 
 an1 · · · an,j−1 bn an,j+1 · · · ann 

is obtained by replacing the jth column of |A| by the column whose components are
b1 , b2 , . . . , bn . If the right-hand side of the equation system (∗) consists only of zeros, so
that it can be written in matrix form as Ax = 0, the system is called homogeneous. A
homogeneous system will always have the trivial solution x1 = x2 = · · · = xn = 0.

33
Lemma 4.2 Ax = 0 has nontrivial solutions if and only if |A| = 0.
I omit the proof of Lemma 4.2.

Exercise 4.9 Use Cramer’s rule to solve the following system of equations:

2x1 − 3x2 = 2
4x1 − 6x2 + x3 = 7
x1 + 10x2 = 1.

4.3 Vectors
An n-vector is an ordered n-tuple of numbers. It is often convenient to regard the rows
and columns of a matrix as vectors, and an n-vector can be understood either as a 1 × n
matrix a = (a1 , a2 , . . . , an ) (a row vector ) or as an n × 1 matrix aT = (a1 , a2 , . . . , an )T
(a column vector). The operations of addition, subtraction and multiplication by scalars
of vectors are defined in the obvious way. The dot product (or inner product) of the
n-vectors a = (a1 , a2 , . . . , an ) and b = (b1 , b2 , . . . , bn ) is defined as
n

a · b = a1 b1 + a2 b2 + · · · + an bn = ai bi
i=1

Proposition 4.3 If a, b, and c are n-vectors and α is a scalar, then

1. a · b = b · a,

2. a · (b + c) = a · b + a · c,

3. (αa) · b = a · (αb) = α(a · b).

4. a · a = 0 =⇒ a = 0

5. (a + b) · (a + b) = a · a + 2(a · b) + b · b.

Exercise 4.10 Prove Proposition 4.1. If you find it difficult to do so, focus on vectors
in R2 .
The Euclidean norm or length of the vector a = (a1 , a2 , . . . , an ) is
√ 
a = a · a = a21 + a22 + · · · + a2n

Note that αa = |α|a for all scalars and vectors.

Lemma 4.3 The following useful inequalities hold.

1. |a · b| ≤ a · b (Cauchy-Schwartz inequality)

34
 2. a + b ≤ a + b (Minkowski inequality)
Proof of Cauchy-Schwartz inequality: Define f (t) as

f (t) = (ta + b) · (ta + b),

where t ∈ R. Because of the definition of dot products, we have f (t) ≥ 0 for any t ∈ R.

f (t) = t2 a2 + 2t(a · b) + b2 .

Then, using the formula, we solve the above equation with respect to t:

−(a · b) ± (a · b)2 − a2 b2
t=
a2
Since f (t) ≥ 0 for any t ∈ R, we must have

|a · b| ≤ ab. 

Exercise 4.11 Prove property 2 in Lemma 4.3. (Hint: It suffices to show that a+b2 ≤
(a + b)2 )
Cauchy-Schwartz inequality implies that, for any a, b ∈ Rn ,
a·b
−1 ≤ ≤ 1.
ab
Thus, the angle θ between nonzero vectors a and b ∈ Rn is defined by
a·b
cos θ = , θ ∈ [0, π]
a · b
This definition reveals that cos θ = 0 if and only if a · b = 0. Then θ = π/2. In symbols,

a⊥b ⇐⇒ a · b = 0

The hyperplane in Rn that passes through the point a = (a1 , . . . , an ) and is orthogo-
nal to the nonzero vector p = (p1 , . . . , pn ), is the set of all points x = (x1 , . . . , xn ) such
that

p · (x − a) = 0

4.4 Linear Independence

Definition 4.1 The n vectors a1 , a2 , . . . , an in Rm are linearly dependent if there
exist numbers c1 , c2 , . . . , cn , not all zero, such that

c1 a1 + c2 a2 + · · · + cn an = 0

If this equation holds only when c1 = c2 = · · · = cn = 0, then the vectors are linearly
independent.

35
Exercise 4.12 a1 = (1, 2), a2 = (1, 1), and a3 = (5, 1) ∈ R2 . Show that a1 , a2 , a3 are
linearly dependent.
Leta1 , a2 , . . . , an ∈ Rn \{0}. Suppose that, for any i = 1, . . . , n, it follows that
ai
= n
j=i λj aj for any λ1 , . . . , λi−1 , λi+1 , . . . , λn ∈ R. Then, the entire space R is
spanned by the set of all linear combinations of a1 , . . . , an .

Lemma 4.4 A set of n vectors a1 , a2 , . . . , an in Rm is linearly dependent if and only if

at least one of them can be written as a linear combination of the others. Or equivalently:
A set of vectors a1 , a2 , . . . , an in Rm is linearly independent if and only if none of them
can be written as a linear combination of the others.
Proof of Lemma 4.4: Suppose that a1 , a2 , . . . , an are linearly dependent. Then
the equation c1 a1 + · · · + cn an = 0 holds with at least one of the coefficients ci differently
from 0. We can, without loss of generality, assume that c1
= 0. Solving the equation for
a1 yields
c2 cn
a1 = − a2 − · · · − an .
c1 c1
Thus, a1 is a linear combination of the other vectors. 

4.4.1 Linear Dependence and Systems of Linear Equations

Consider the general system of m equations in n unknowns.
a11 x1 + a12 x2 + · · · + a1n xn = b1
a21 x1 + a22 x2 + · · · + a2n xn = b2 ⇐⇒ x1 a1 + · · · + xn an = b (∗)
··························· ··· ···
an1 x1 + an2 x2 + · · · + amn xn = bm
Here a1 , . . . , an are the column vectors of coefficients, and b is the column vector with
components b1 , . . . , bm .

Suppose that (∗) has two solutions (u1 , . . . , un ) and (v1 , . . . , vn ). Then,
u1 a1 + · · · + un an = b and v1 a1 + · · · + vn an = b
Subtracting the second equation from the first yields
(u1 − v1 )a1 + · · · + (un − vn )an = 0.
Let c1 = u1 −v1 , . . . , cn = un −vn . The two solutions are different if and only if c1 , . . . , cn
are not all equal to 0. We conclude that if system (∗) has more than one solution, then
the column vectors a1 , . . . , an are linearly dependent. 2 Equivalently, If the column
vectors a1 , . . . , an are linearly independent, then system (∗) has “at most” one solution.
3

2
Recall Lemma 4.4.
3
Is there anything to say when there is no solution? The answer is yes. I can use Farkas Lemma to
check if there is any solution to the system. See Appendix 1 in this chapter for Farkas Lemma.

36
Theorem 4.1 The n column vectors a1 , a2 , . . . , an of the n × n matrix
⎛ ⎞ ⎛ ⎞
a11 a12 · · · a1n a1j
⎜ a21 a22 · · · a2n ⎟ ⎜ a2j ⎟
⎜ ⎟ ⎜ ⎟
A=⎜ . .. .. .. ⎟ , where aj = ⎜ .. ⎟ j = 1, . . . , n
⎝ .. . . . ⎠ ⎝ . ⎠
an1 an2 · · · ann anj

are linearly independent iff |A|

= 0.
Proof of Theorem 4.1: The vectors a1 , . . . , an are linearly independent iff the
vector equation x1 a1 + · · · xn an = 0 has only the trivial solution x1 = · · · = xn = 0. This
vector equation is equivalent to a homogeneous systems of equations, and therefore, it
has only the trivial solution iff |A|
= 0. 4 

Definition 4.2 The rank of a matrix A, written r(A), is the maximum number of
linearly independent column vectors in A. If A is the 0 matrix, we put r(A) = 0.

4.5 Eigenvalues
4.5.1 Motivations
Consider the matrix A below.


2 0
A=
0 3

Let x, y be vectors in R2 . Suppose that y is derived from Ax as follows.





y1 2 0 x1 2x1
= = ⇐⇒ y = Ax
y2 0 3 x2 3x2

The linear transformation (matrix) A extends x1 into 2x1 along the x1 axis and x2
into 3x2 along the x2 axis. Importantly, there is no interaction between x1 and x2 through
the linear transformation A. This is, I believe, a straightforward extension of the linear
transformation in R into Rn . Define e1 = (1, 0) and e2 = (0, 1) as the unit vectors in R2 .
Then, x = x1 e1 + x2 e2 and y = 2x1 e1 + 3x2 e2 . In other words, (e1 , e2 ) is the unit vector
in the original space and (2e1 , 3e2 ) is the unit vector in the space transformed through
A. Next, consider the matrix B as follows.


1 1
B=
−2 4

Now, we don’t have a clear image about what is going on through the linear trans-
formation A. However, consider the following different unit vectors f1 = (1, 1) and
4
Check Section 4.2.3 for this argument. Recall that the system of linear equations is homogeneous if
it is expressed by Ax = 0.

37
f2 = (1, 2). Then,
f1 f2


  


  

1 1 1 1 1 1 1 1
=2 , and =3
−2 4 1 1 −2 4 2 2
           
A 2f1 A 3f2

This shows that once we take f1 and f2 as the new coordinate system, the linear
transformation B is the same as A but now along the f1 and f2 axes, respectively.
Finally, consider the matrix C below.


2 −3
C=
4 2
It turns out that there is no way of finding the new coordinate system in which
the linear transformation C can be seen as either extending or shrinking the vectors in
each new axis. The reason why we don’t find such a new coordinate system is that we
restrict our attention to Rn . Once we allow for the unit vectors in the new system to be
complex numbers, we again will be successful to find the new coordinate system in which
everything is easy to understand. 5

Consider another different unit vector

 √   √ 
2 3 2 3
f˜1 = 1, i and f˜2 = 1, − i .
3 3

Then,
f˜1
  

  
2 −3 1 √ 1
√ = (2 − 2 3i) √ , and
4 2 2 3i 2 3i
   3 3
  
A √
(2−2 3i)f˜1
f˜2
  

  
2 −3 1√ √ 1√
= (2 + 2 3i)
4 2 − 2 33i − 2 33i
     
A √
(2+2 3i)f˜2

I want to generalize the above argument. Suppose there happens to be a scalar λ

with the special property that
Ax = λx (∗)
In this case, we would have A2 x = A(Ax) = A(λx) = λAx = λλx = λ2 x, in general,
An x = λn x.
5
Those who are interested in the definition of complex numbers should be referred to Appendix 2 in
this chapter.

38
Definition 4.3 If A is an n × n matrix, then a scalar λ is an eigenvalue of A if there
is a nonzero vector x ∈ Rn such that

Ax = λx

Then x is an eigenvector of A (associated with λ)

4.5.2 How to Find Eigenvalues

The eigenvalue equation can be written as

(A − λI)x = 0

where I denotes the identity matrix of order n. Note that this linear system of equations
has a solution x
= 0 if and only if the coefficient matrix has determinant equal to 0 –
that is, iff |A − λI| = 0. Letting p(λ) = |A − λI|, where A = (aij )n×n , we have the
equation
 
 a11 − λ a12 ··· a1n 

 a21 a22 − λ · · · a2n 

p(λ) = |A − λI| =  .. .. . ..  = 0.
 . . .. . 
 
 an1 an2 · · · ann − λ 

This is called the characteristic equation of A. From the definition of determinant,

it follows that p(λ) is a polynomial of degree n in λ. According to the fundamental
theorem of algebra, it has exactly n roots (real or complex), provided that any multiple
roots are counted appropriately.

Theorem 4.2 (The Fundamental Theorem of Algebra by Gauss (1779)) Consider

a polynomial of degree n in z shown below.

z n + an−1 z n−1 + · · · + a1 z + a0 = 0, (∗)

where a0 , . . . an−1 ∈ C. Then, (∗) has n solutions z1∗ , . . . , zn∗ with the property that zi∗ ∈ C
for each i = 1, . . . , n. This includes the case in which zi∗ = zj∗ for some i
= j.

Exercise 4.13 Find the eigenvalues and the associated eigenvectors of the matrices, A
and B.


1 2
A =
3 0


0 1
B =
−1 0

In fact, it is convenient to write the characteristic function as a polynomial in −λ:

p(λ) = (−λ)n + bn−1 (−λ)n−1 + · · · + b1 (−λ) + b0

39
The zeros of this characteristic polynomial are precisely the eigenvalues of A. Denoting
the eigenvalues by λ1 , λ2 , . . . , λn ∈ C, we have

p(λ) = (−1)n (λ − λ1 )(λ − λ2 ) · · · (λ − λn )

Theorem 4.3 If A is an n × n matrix with eigenvalues λ1 , λ2 , . . . , λn , then

1. |A| = λ1 λ2 · · · λn

2. tr(A) = a11 + a22 + · · · + ann = λ1 + λ2 + · · · + λn

Proof of Theorem 4.3: Putting λ = 0, we see that p(0) = b0 = |A|. Specifically,
λ = 0 gives p(0) = (−1)n (−1)n λ1 λ2 · · · λn . Since (−1)n (−1)n = ((−1)n )2 = 1, we have
b0 = |A| = λ1 λ2 · · · λn . The product of the elements on the main diagonal of |A − λI| is

(a11 − λ)(a22 − λ) · · · (ann − λ).

If we choose ajj from one of these parentheses and −λ from the remaining n − 1, then
add over j = 1, . . . n, we obtain the term

(a11 + a22 + · · · + ann )(−λ)n−1

Since we cannot obtain other terms with (−λ)n−1 except the above terms, we conclude
that bn−1 = a11 + a22 + · · · ann , the trace of A. 

4.6 Diagonalization
Let A and P be n × n matrices with P invertible. Then A and P −1 AP have the same
eigenvalues. This is true because the two matrices have the same characteristic polyno-
mial:

|P −1 AP − λI| = |P −1 AP − P −1 λIP | = |P −1 (A − λI)P |

= |P −1 ||A − λI||P | = |A − λI|

where we use the fact that |P −1 | = 1/|P | and |AB| = |A||B| (See Proposition 4.1 and
4.2.) .

An n × n matrix A is diagonalizable if there exists an invertible n × n matrix P and

a diagonal matrix D such that P −1 AP = D.

Theorem 4.4 (Diagonalization Theorem) An n × n matrix A is diagonalizable if

and only if it has a set of n linearly independent eigenvectors x1 , . . . , xn ∈ Cn . In this
case,

P −1 AP = diag(λ1 , . . . , λn )

where P is the matrix with x1 , . . . , xn ∈ Cn as its columns, and λ1 , . . . , λn ∈ C are the

corresponding eigenvalues.

40
 Proof of Diagonalization Theorem: (⇐=) Suppose that A has n linearly indepen-
dent eigenvectors x1 , . . . , xn , with corresponding eigenvalues λ1 , . . . , λn . Let P denote
the matrix whose columns are x1 , . . . , xn . Then, AP = P D (⇔ AP = DP because D is
diagonal), where D = diag(λ1 , . . . , λn ). Because the eigenvectors are linearly indepen-
dent, P is invertible, so P −1 AP = D. (=⇒) If A is diagonalizable, there exists invertible
n × n matrix P such that P −1 AP = D. Then, AP = P D. Since D is a diagonal matrix
by our hypothesis, we have AP = DP . The columns of P must be eigenvectors of A,
and the diagonal elements of D must be the corresponding eigenvalues. 

A matrix P is said to be orthogonal if P T = P −1 , i.e., P T P = I. If x1 , . . . , xn are the

 
n column vectors of P , then x1 , . . . , xn are the row vectors of the transformed matrix,

P . The condition P T P = I then reduces to the n2 equations xTi · xj = 1 if i = j and
xTi · xj = 0 if i
= j.

Theorem 4.5 If the matrix A = (aij )n×n is symmetric, then:

1. All the n eigenvalues λ1 , . . . , λn are real.

2. Eigenvectors that corresponds to different eigenvalues are orthogonal.

3. There exists an orthogonal matrix P (i.e., P T = P −1 ) such that

⎛ ⎞
λ1 0 · · · 0
⎜ 0 λ2 · · · 0 ⎟
⎜ ⎟
P −1 AP = ⎜ . .. . . .. ⎟
⎝ .. . . . ⎠
0 0 ··· λn

The columns v1 , v2 , . . . , vn of the matrix P are eigenvectors of unit length corre-

sponding to the eigenvalues λ1 , λ2 , . . . , λn .

Proof of Theorem 4.5: (1) We will show this for n = 2. The eigenvalues of 2 × 2
matrix A are given by the quadratic equation.
 
 a11 − λ a 
|A − λI| =  12  = λ2 − (a11 + a22 )λ + (a11 a22 − a12 a21 ) = 0 (∗)
a21 a22 − λ 

The roots of the quadratic equation (∗) are


(a11 + a22 ) ± (a11 + a22 )2 − 4(a11 a22 − a12 a21 )
λ= .
2
These roots are real if and only if (a11 + a22 )2 ≥ 4(a11 a22 − a12 a21 ), which is equivalent
to

(a11 − a22 )2 + 4a12 a21 ≥ 0

(a11 − a22 )2 + 4a212 ≥ 0 if a12 = a21 because of the symmetry of A

41
This is indeed the case. (2) Suppose that Axi = λi xi and Axj = λj xj by λi
= λj .
Multiplying these equalities from the left by xTj and xTi , respectively,

xTj Axi = λi xj xi and xTi Axj = λj xTi xj

Applying transpose operation on both hand sides, we obtain

xTi AT xj = λi xTi xj and xTj AT xi = λj xTj xi

Since A is symmetric so that A = AT , the above is simplified to

xTi Axj = λi xTi xj and xTj Axi = λj xTj xi


This implies that (λi − λj )xi xj = 0. Since λi − λj
= 0 by our hypothesis, we must

have xi xj = 0 and thus, xi and xj are orthogonal. (3) Suppose all the real – which
we are supposed to know from (1) – eigenvalues are different. 6 Then, according to (2)
which we have shown above, the associated eigenvectors are mutually orthogonal. Hence,
the eigenvectors are linearly independent. When we define P as the collection of the
eigenvectors x1 , . . . xn as its columns, it follows that P −1 = P T . By the diagonalization
theorem (Theorem 4.4), A is diagonalizable. We can choose the eigenvectors so that they
all have length 1, by replacing each xj with xj /xj . 

Exercise 4.14 Let a 2 × 2 symmetric matrix A be given below.


2 1
A=
1 2
Compute the matrix P described in Theorem 4.5.

4.7 Quadratic Forms

A quadratic form in n variables is a function Q of the form
n
n
Q(x1 , . . . , xn ) = aij xi xj = a11 x211 + a12 x1 x2 + · · · + aij xi xj + · · · + ann x2n .
i=1 j=1

where the aij are constants. Suppose we put x = (x1 , . . . , xn )T and A = (aij ). Then, it
follows from the definition of matrix multiplication that

Q(x1 , . . . , xn ) = Q(x) = xT Ax.

Of course, xi xj = xj xi , so we can write aij xi xj + aji xj xi = (aij + aji )xi xj . If we replace

aij and aji by (aij + aji )/2, then the new numbers aij and aji become equal without
changing Q(x). Thus, we can assume that aij = aji for all i and j, which means that
the matrix A is symmetric. Then A is called the symmetric matrix associated with Q,
and Q is called a symmetric quadratic form.
6
We omit the case where some of the eigenvalues are equal.

42
 
Definition 4.4 A quadratic form Q(x) = x Ax, as well as its associated symmetric
matrix A, are said to be positive definite, positive semidefinite, negative definite,
or negative semidefinite according as

Q(x) > 0, Q(x) ≥ 0, Q(x) < 0, Q(x) ≤ 0,

for all x ∈ R\{0}. The quadratic form Q(x) is indefinite if there exist vectors x∗ and
y ∗ such that Q(x∗ ) < 0 and Q(y ∗ ) > 0.
Let A = (aij ) be any n × n matrix. An arbitrary principal minor of order r is the
determinant of the matrix obtained by deleting all but r rows and r columns in A with
the same numbers. In particular, a principal minor of order r always includes exactly r
elements of the main (principal) diagonal. We call the determinant |A| itself a principal
minor (no rows and columns are deleted). A principal minor is said to be a leading
principal minor of order r (1 ≤ r ≤ n), if it consists of the first “leading” rows and
columns of |A|.

Suppose A is an arbitrary n × n matrix. The leading principal minors of A are

 
 a11 a12 · · · a1k 

 a21 a22 · · · a2k 

Dk =  . .. .. ..  , k = 1, . . . , n
 .. . . . 

 ak1 ak2 · · · akk 

Exercise 4.15 Consider a 3 × 3 matrix A:

⎛ ⎞
a11 a12 a13
A = ⎝ a21 a22 a23 ⎠ .
a31 a32 a33

Compute all the principal minors of A.

Theorem 4.6 Consider the quadratic form

n
n
Q(x) = aij xi xj (aij = aji )
i=1 j=1

with the associated symmetric matrix A = (aij )n×n . Let Dk be the leading principal
minor of A of order k and let Δk denote an arbitrary principal minor of order k. Then
we have

1. Q is positive definite ⇐⇒ Dk > 0 for k = 1, . . . , n

2. Q is positive semidefinite ⇐⇒ Δk ≥ 0 for all principal minors of order k =

1, . . . , n.

3. Q is negative definite ⇐⇒ (−1)k Dk > 0 for k = 1, . . . , n

43
 4. Q is negative semidefinite ⇐⇒ (−1)k Δk ≥ 0 for all principal minors of order
k = 1, . . . , n.

Proof of Theorem 4.6: We only prove this for n = 2. Then, the quadratic form is

Q(x1 , x2 ) = a11 x21 + 2a12 x1 x2 + a22 x22

After some manipulation through perfect square, we obtain

2

a12 a212
Q(x1 , x2 ) = a11 x1 + x2 + a22 − x22
a11 a11 
   ≥0
>0

Thus, we obtain

• Q(x1 , x2 ) > 0 ⇐⇒ a11 > 0 and a11 a22 − a212 > 0.

• Q(x1 , x2 ) < 0 ⇐⇒ a11 < 0 and a11 a22 − a212 < 0. 


Theorem 4.7 Let Q = x Ax be a quadratic form, where the matrix A is symmetric,
and let λ1 , . . . , λn be the (real) eigenvalues of A. Then,

1. Q is positive definite ⇐⇒ λ1 > 0, . . . , λn > 0

2. Q is positive semidefinite ⇐⇒ λ1 ≥ 0, . . . , λn ≥ 0

3. Q is negative definite ⇐⇒ λ1 < 0, . . . , λn < 0

4. Q is negative semidefinite ⇐⇒ λ1 ≤ 0, . . . , λn ≤ 0

5. Q is indefinite ⇐⇒ A has eigenvalues with opposite signs.

Proof of Theorem 4.7: According to Theorem 4.5, there exists an orthogonal


matrix P such that P AP = diag(λ1 , . . . , λn ). Let y = (y1 , . . . , yn )T be the n × 1 matrix
defined by y = P T x. Then, x = P y, so that

xT Ax = (P y)T AP y = y T P T AP y = y T diag(λ1 , . . . , λn )y = λ1 y12 + λ2 y22 + · · · + λn yn2

Also, x = 0 iff y = 0. This completes the proof. 

4.8 Appendix 1: Farkas Lemma

In this Appendix, I follow “Advanced Mathematical Economics,” by Rakesh Vohra.

44
4.8.1 Preliminaries
Definition 4.5 A vector y can be expressed as a linear combination of a vectors in
S = {x1 , x2 , . . . } if there are real numbers {λj }j∈S such that

y= λj xj
j∈S

The set of all vectors that can be expressed as a linear combination of vectors in S is
called the span of S and denoted span(S).

Definition 4.6 The rank of a (not necessarily finite) set S of vectors is the size of the
largest subset of linearly independent vectors in S.

Definition 4.7 Let S be a set of vectors and B ⊂ S be finite and linearly independent.
The set B of vectors is said to be a maximal linear independent set if the set B ∪ {x}
is linearly dependent for all vectors x ∈ S\B. A maximal linearly independent subset of
S is called a basis of S.

Theorem 4.8 Every S ⊂ Rn has a basis. If B is a basis of S, then span(S) = span(B).

 
Theorem 4.9 Let S ⊂ Rn . If B and B are two bases of S, then |B| = |B |.
From this theorem, one can see that if S has a basis B, then the rank of S and B.

Definition 4.8 Let S be a set of vectors. The dimension of span(S) is the rank of S.

Definition 4.9 The kernel or null space of A is the set {x ∈ Rn |Ax = 0}.
The following theorem summarizes the relationship between the span of A and its
kernel.

Theorem 4.10 If A is an m × n matrix, then the dimension of span(A) plus the dimen-
sion of the kernel of A is n.
This is sometimes written as

dim[span(A)] + dim[ker(A)] = rank(A) + dim[ker(A)] = n.

The column rank of a matrix is the dimension of the span of its columns. Similarly,
the row rank is the dimension of the span of its row.

Theorem 4.11 Let A be an m × n matrix. Then, the column rank of A and AT (the
tranpose of A) are the same.
Thus, the column and row rank of A are equal. This allows us to define the rank of
a matrix A to be the dimension of span(A).

45
4.8.2 Fundamental Theorem of Linear Algebra
Let A be an m × n matrix of real numbers. We will be interested in problems of the
following kind:

Given b ∈ Rm , find an x ∈ Rn such that Ax = b or prove that no such x exists.

Convincing another that Ax = b has a solution (when it does) is easy. One merely
exhibits and they can verify that the solution does indeed satisfy the equations. What
if the system Ax = b does not admit a solution? By framing the problem in the right
way, we can bring to bear the machinery of linear algebra. Specifically, given b ∈ Rm ,
the problem of finding an x ∈ Rn such that Ax = b can be stated as: is b ∈ span(A)?

Theorem 4.12 (Gauss (??)) Let A be an m × n matrix, b ∈ Rm and F = {x ∈

Rn |Ax = b}. Then, either F
= ∅ or there exists y ∈ Rm such that yA = 0 and yb
= 0
but not both.
Suppose F = ∅. Then, b is not in the span of the columns of A. If I think of the space
of the columns of A as a plane, then b is a vector pointing out of the plane. Thus, any
vector y orthogonal to this plane (and so to every column of A) must have a non-zero
dot product with b. Now for an algebraic interpretation. Take any linear combination
of the equations in the system Ax = b. This linear combination can be obtained by
pre-multiplying each side of the equation by a suitable vector y, i.e., yAx = yb. Suppose
there is a solution x∗ to the system, i.e., Ax∗ = b. Any linear combination of these
equations results in an equation that x∗ satisfies as well. In particular, x∗ must also be
a solution to the resulting equation: yAx = yb. Suppose I found a vector y such that
yAx
= yb then clearly the original system Ax = b could not have a solution.

Proof: First, we prove the “not both” part. Suppose that F

= ∅. Choose any x ∈ F .
Then, for any y ∈ Rm , we have

yb = yAx = (yA)x = 0,

which contradicts the fact that yb

= 0 for some y.
If F
= ∅, we are done. Suppose that F
= ∅. Hence, b cannot be in the span of the

columns of A. Thus, the rank of C = [A, b], r , is one larger than the rank, r, of A. That

is, r = r + 1. Since C is a m × (n + 1) matrix,

rank(C T ) + dim[ker(C T )] = m = rank(AT ) + dim[ker(AT )].

Using the fact that the rank of a matric and its transpose coincide, we have

r + dim[ker(C T )] = m = r + dim[ker(AT )],

i.e., dim[ker(C T )] = dim[ker(AT )] − 1. Since the dimension of ker(C T ) is one smaller

than the dimension of ker(AT ), we can find a y ∈ ker(AT ) that is not in ker(C T ). Hence,
yA = 0 but yb
= 0. 

46
4.8.3 Linear Inequalities
Now consider the following problem:

Given a b ∈ Rm , find an x ∈ Rn such that Ax ≤ b or show that no such x exists.

The problem differs from the earlier one in that “=” has been replaced by “≤.”

4.8.4 Non-Negative Solutions

I focus on finding a non-negative x ∈ Rn such that Ax = b or show that no such x exists.
Observe that if b = 0, the problem is trivial, so I assume that b
= 0.

Definition 4.10 A set C of vectors is called a cone if λx ∈ C whenever x ∈ C and

λ > 0.

Definition 4.11 The set of all non-negative linear combinations of the columns of A is
called the finite cone generated by the columns of A. It is denoted cone(A).
Note the difference between span(A) and cone(A) below:

span(A) = {y ∈ Rm |y = Ax for some x ∈ Rn }

and

cone(A) = y ∈ Rm |y = Ax for some x ∈ Rn+

Theorem 4.13 (Farkas Lemma (1902)) Let A be an m × n matrix, b ∈ Rm and

F = {x ∈ Rn |Ax = b, x ≥ 0}. Then, either F
= ∅ or there exists y ∈ Rm such that
yA ≥ 0 and y · b < 0 but not both.
Take any linear combination of the equations in Ax = b to get yAx = yb. A non-
negative solution to the first system is a solution to the second. If we can choose y so that
yA ≥ 0 and y · b < 0, we find that the left hand side of the single equation yAx = yb is
at least zero while the right hand side is negative, a contradiction. Thus, the fist system
cannot have a non-negative solution.

Proof: First we prove that both statements cannot hold simultaneously. Suppose
not. Let x∗ ≥ 0 be a solution to Ax = b and y ∗ a solution to yA ≥ 0 such that
y ∗ b < 0. Notice that x∗ must be a solution to y ∗ Ax = y ∗ b. Thus, y ∗ Ax∗ = y ∗ b. The
0 ≤ y ∗ Ax∗ = y ∗ b < 0, which is a contradiction.

If b ∈
/ span(A) (i.e., there is no x such that Ax = b), by the previous theorem
(Theorem ??), there is a y ∈ Rm such that yA = 0 and yb
= 0. If it so happens that the
given y has the property that yb < 0 we are done. If yb > 0, then negate y and again we
are done. So, we may suppose that b ∈ span(A) but b ∈ / cone(A), i.e., F = ∅.

47
 Let r be the rank of A. Note that n ≥ r. Since A contains r linearly independent
column vectors and b ∈ span(A), we can express b as a linear combination of
an r-subset
D of linearly independent columns of A. Let D = {a , . . . , a } and b = rt=1 λit ait .
i1 ir

Note that D is linearly independent. Since b ∈ / cone(A), at least one of {λit }t≥1 is
negative.

Now apply the following four step procedure repeatedly. Subsequently, we show that
the procedure must terminate.
1. Choose the smallest index h amongst {i1 , . . . , ir } with λh < 0.
2. Choose y so that y · a = 0 for all a ∈ D\ah and y · ah
= 0. This can be done by the
previous theorem (Theorem ??) because ah ∈ / span(D\ah ). Normalize y so that
h
y · a = 1. Observe that y · b = λh < 0.
3. If y · aj ≥ 0 for all columns aj of A stop, and the proof is complete.
4. Otherwise, choose the smallest index w amongst {1, . . . , n} such that y · aw < 0.
Note that aw ∈
/ D\ah . Replace D by {D\ah } ∪ aw , i.e., exchange ah for aw .
To complete the proof, we must show that the procedure terminates. Let D k denote
the set D at the start of the kth iteration of the four step procedure described above. If
the procedure does not terminate, there is a pair (k,
) with k <
such that D k = D ,
i.e., the procedure cycles.

Let s be the largest index for which as has been removed from D at the end of one
of the iterations k, k + 1, . . . ,
− 1, say p. Since D = Dk , there is a q such that as is
inserted into D q at the end of iteration q, where k ≤ q <
. No assumption is made
about whether p < q or p > q. Notice that
Dp ∩ {as+1 , . . . , an } = Dq ∩ {as+1 , . . . , an }.

Let D p = {ai1 , . . . , air }, b = λi1 ai1 + · · · λir air and let y be the vector found in step two
of iteration q. Then:
  !  
0 > y · b = y λi1 ai1 + · · · + λir air = y λi1 ai1 + · · · + y λir air > 0,
which is a contradiction. The first inequality comes from the previous theorem (Theorem
??). To see why the last inequality must be true:
• When ij < s, we have from Step 1 of iteration p that λij ≥ 0. From Step 4 of

iteration q, we have y · aij ≥ 0.
• When ij = s, we have from Step 1 of iteration p that λij < 0. From Step 4 of

iteration q we have y · aij < 0.
• When ij > s, we have from Dp ∩ {as+1 , . . . , an } = Dq ∩ {ar+1 , . . . , an } and Step 2

of iteration q that y · aij = 0.
This complete the proof. 

48
4.8.5 The General Case
The problem of deciding whether the system {x ∈ Rn | Ax ≤ b} has a solution can be
reduced to the problem of deciding if Bz = b, z ≥ 0 has a solution for a suitable matrix
B. 
First observe that any inequality of the form j aij xj ≥ bi can be turned into an
equation by the subtraction of a surplus variable, s. That is, define a new variable si ≥ 0
such that

aij xj − si = bi .
j

Similarly, an inequality of the form j aij xj ≤ bi can be converted into an equation by
the addition of a slack variable, si ≥ 0 as follows:

aij xj + si = bi .
j

A variable xj that is unrestricted in sign can be replaced by two non-negative variables

 
zj and zj by setting xj = zj − zj . In this way any inequality system can be converted
into an equality system with non-negative variables. We will refer this as converting into
standard form.

As an example, we derive the Farkas alternative for the system {x|Ax ≤ b, x ≥ 0}.
Deciding solvability of Ax ≤ b for x ≥ 0 is equivalent to solvability of Ax + Is = b where
x, s ≥ 0. Set B = [A|I] and z = (x, s)T and we can write the system as Bz = b, z ≥ 0.
Now apply the Farkas lemma to this system:

yB ≥ 0, yb < 0.

Now 0 ≤ yB = y[A|I] implies yA ≥ 0 and y ≥ 0. So, the Farkas alternative is {y|yA ≥

0, y ≥ 0, yb < 0}. The principle here is that by a judicious use of auxiliary variable,
one can convert almost anything into standard form.

4.9 Appendix 2: Linear Spaces

4.9.1 Number Fields
Linear algebra makes use of number systems (number fields). By a number field I mean
any set K of objects, called “numbers,” which, when subjected to four arithmetic op-
erations again give elements of K. More exactly, these operations have the following
properties F 1, F 2, and F 3 (field axioms):

F1: To every pair of numbers α and β in K, there corresponds a (unique) number

α + β in K, called the sum of α and β, where

1. α + β = β + α ∀α, β ∈ K (addition is communicative);

49
 2. (α + β) + γ = α + (β + γ) ∀α, β, γ ∈ K (addition is associative);

3. There exists a number 0 (zero) in K such that 0 + α = α ∀α ∈ K;

4. For every α ∈ K, there exists a number (negative element) γ ∈ K such that

α + γ = 0.

The solvability of the equation α + γ = 0 ∀α allows us to carry out the operation of

subtraction, by defining the difference β − α as the sum of the number β and the solution
γ of the equation α + γ = 0.

F2: To every pair of numbers α, β ∈ K, there corresponds a (unique) number α · β (or

αβ) in K, called the product of α and β, where

1. αβ = βα ∀α, β ∈ K (multiplication is commutative);

2. (αβ)γ = α(βγ) ∀α, β, γ ∈ K (multiplication is associative);

3. There exists a number 1 (

= 0) in K such that 1 · α = α ∀αinK;

4. For every α
= 0 ∈ K, there exists a number (reciprocal element) γ ∈ K such that
αγ = 1.

F3: Multiplication is distributive over addition, i.e., for every α, β, γ ∈ K,

α(β + γ) = αβ + αγ.

The solvability of the equation αγ = 1 for every α

= 0 allows us to carry out the
operation of division, by defining the quotient β/α as the product of the number β and
the solution γ of the equation αγ = 1.

The numbers 1, 1 + 1 = 2, 2 + 1 = 3, etc. are said to be natural; it is assumed

that none of these numbers is zero. The set of natural numbers is denoted as N. By
the integers in a field K, we mean the set of all natural numbers together with their
negatives and the number zero. The set of integers is denoted as Z. By the rational
numbers in a field K, we mean the set of all quotients p/q, where p and q are integers
and q
= 0. The set of rational numbers is denoted as Q.

Two field K and K are said to be isomorphic if we can set up a one-to-one cor-

respondence between K and K such that the number associated with every sum (or

product) of numbers in K is the sum (or product) of the corresponding numbers in K .
The number associated with every difference (or quotient) of numbers in K will then be

the difference (or quotient) of the corresponding numbers in K .

The most commonly encountered concrete examples of number fields are the follow-
ing:

50
 1. The field of rational numbers, i.e., of quotients p/q where p and q
= 0 are the
ordinary integers subject to the ordinary operations of arithmetic. It should be
noted that the integers by themselves do not form a field, since they do not satisfy
axiom F2-4. It follows that every field K has a subset isomorphic to the field of
rational numbers.

2. The field of real numbers, having the set of all points of the real line as its geometric
counterpart. The set of real numbers is denoted as R. An axiomatic treatment of
the field of real numbers is achieved by supplementing axioms F 1, F 2, F 3 with the
axioms of order and the least upper bound principle.

3. The field of complex numbers of the form a + ib, where a and b are real numbers
(i is not a real number), equipped with the following operations of addition and
multiplication:

(a1 + ib1 ) + (a2 + ib2 ) = (a1 + a2 ) + i(b1 + b2 ),

(a1 + ib1 )(a2 + ib2 ) = (a1 a2 − b1 b2 ) + i(a1 b2 + a2 b1 ).

The set of complex numbers is denoted as C. For numbers of the form a + i0, these
operations reduce to the corresponding operations for real numbers; briefly I write
a + i0 = a and call complex numbers of this form real. Thus, it can be said that
the field of complex numbers has a subset isomorphic to the field of real numbers.
Complex numbers of the form 0 + ib are said to be (purely) imaginary and are
designated briefly by ib. It follows from the multiplication rule that

i2 = i · i = (0 + i1)(0 + i1) = −1.

4.9.2 Definitions
The concept of a linear space generalizes that of the set of all vectors. The generalization
consists first in getting away from the concrete nature of the objects involved (directed
line segments) without changing the properties of the operations on the objects, and
secondly in getting away from the concrete nature of the admissible numerical factors
(real numbers). This leads the following definition.

Definition 4.12 A set V is called a linear (or affine) space over a field K if

1. Given any two elements x, y ∈ V , there is a rule (the addition rule) leading to a
(unique) element x + y ∈ V , called the sum of x and y;

2. Given any element x ∈ V and any number λ ∈ V , there is a rule (the multiplication
by a number) leading to a (unique) element λx ∈ V , called the product of the
element x and the number λ;

3. These two rules obey the axioms listed below, VS1 and VS2.

VS 1: The addition rule has the following properties:

51
 1. x + y = y + x for every x, y ∈ V ;
2. (x + y) + z = x + (y + z) for every x, y, z ∈ V ;
3. There exists an element 0 ∈ V (the zero vector ) such that x + 0 = x for every
x∈V;
4. For every x ∈ V , there exists an element y ∈ V (the negative element) such that
x + y = 0.

VS 2: The rule for multiplication by a number has the following properties:

1. 1 · x = x for every x ∈ V ;
2. α(βx) = (αβ)x for every x ∈ V and α, β ∈ K;
3. (α + β)x = αx + βx for every x ∈ V and α, β ∈ K;
4. α(x + y) = αx + αy for every x ∈ V and every α ∈ K.

4.9.3 Bases, Components, Dimension

Definition 4.13 A system of linearly independent vectors e1 , e2 , . . . , en in a linear space
V over a field K is called a basis for V if, given any x ∈ V , there exists an expansion

x = ξ1 e1 + ξ2 e2 + · · · + ξn en (∗)

where ξj ∈ V for every j = 1, . . . , n.

It is easy to see that under these conditions, the coefficients in the expansion (∗) are
uniquely determined. In fact, we can write two expansions

x = ξ1 e1 + ξ2 e2 + · · · + ξn en ,
x = η1 e1 + η2 e2 + · · · + ηn en

for a vector x, then, subtracting them term by term, we obtain the relation

0 = (ξ1 − η1 )e1 + (ξ2 − η2 )e2 + · · · + (ξn − ηn )en ,

from which, by the assumption that the vectors e1 , e2 , . . . , en are linearly independent,
we find that

ξ1 = η1 , ξ2 = η2 , . . . , ξn = ηn .

The uniquely defined numbers ξ1 , . . . , ξn are called the components of the vector x with
respect to the basis e1 , . . . , en

The fundamental significance of the concept of a basis for a linear space consists
in the fact that when a basis is specified, the originally abstract linear operations in
the space become ordinary linear operations with numbers, i.e., the components of the
vectors with respect to the given basis. In fact, we have the following.

52
Theorem 4.14 When two vectors of a linear space V are added, their components (with
respect to any basis) are added. When a vector is multiplied by a number λ, all its
components are multiplied by λ.
If, in a linear space V , we can find n linearly independent vectors while every n + 1
vectors of the space are linearly dependent, then the number n is called the dimension
of the space V and the space V itself is called n-dimensional. A linear space in which
we can find an arbitrarily large number of linearly independent vectors is called infinite-
dimensional.

Theorem 4.15 In a space V of dimension n, there exists a basis consisting of n vectors.

Moreover, any set of n linearly independent vectors of the space V is a basis for the space.

Theorem 4.16 If there is a basis in the space V , then the dimension of V equals the
number of basis vectors.

4.9.4 Subspaces
Definition 4.14 (Subspaces) Suppose that a set W of elements of a linear space V
has the following properties:

1. If x, y ∈ W , then x + y ∈ W ;

2. If x ∈ W and λ is an element of the field K, then λx ∈ W .

Then, every set W ⊂ V with properties 1 and 2 above is called linear subspace (or
simply a subspace) of the space V .

Definition 4.15 (The Direct Sum) A linear space W is the direct sum of given
subspaces W1 , . . . , Wm ⊂ W if the following two conditions are satisfied:

1. For every x ∈ W , there exists an expansion

x = x1 + · · · + xm ,

where x1 ∈ W1 , . . . , xm ∈ Wm ;

2. This expansion is unique, i.e., if

x = x1 + · · · + xm = y1 + · · · + ym

where xj , yj ∈ Wj (j = 1, . . . , m), then

x1 = y1 , . . . , xm = ym .

Theorem 4.17 Let W1 be a fixed subspace of an n-dimensional space Vn . Then, there

always exists a subspace W2 ⊂ Vn such that the whole space Vn is the direct sum of W1
and W2 .

53
4.9.5 Morphisms of Linear Spaces
Definition 4.16 Let ϕ be a rule which assigns to every given vector x of linear space V

a vector x in a linear space. Then, ϕ is called morphism (or linear operator) if the
following two conditions hold:

1. ϕ(x + y) = ϕ(x) + ϕ(y) for every x, y ∈ V ;

2. ϕ(αx) = αϕ(x) for every x ∈ V and every α ∈ K.


A morphism ϕ mapping V onto all of V in a one-to-one fashion is called an isomor-

phism, and the spaces V and V themselves are said to be isomorphic (more exactly,
K-isomorphic).

Theorem 4.18 Any two n-dimensional spaces V and V (over the same field K) are
K-isomorphic.

Corollary 4.1 Every n-dimensional linear space over a field K is K-isomorphic to the
space K n . In particular, every n-dimensional complex space is C-isomorphic to the space
Cn , and every n-dimensional real space is R-isomorphic to the space Rn .

54
Chapter 5

Calculus

5.1 Functions of a Single Variable

Roughly speaking, a function y = f (x) is differentiable if it is both continuous and
“smooth,” with no breaks or kinks. The derivative of f is a function giving, at each
value of x, the slope of change in f (x). We sometimes write

dy 
= f (x).
dx

to indicate that f (x) gives us the (instantaneous) amount, dy, by which y changes per
unit change, dx, in x. If the first derivative is a differentiable function, we can take its
derivative which gets the second derivative of the original function

d2 y 
= f (x).
dx2
 
If a function possesses a continuous derivatives f , f , . . . , f n , it is called n-times contin-
uously differentiable, or a C n function. Some rules of differentiation is provided below:

• For constants, α: d/dx(α) = 0.

 
• For sums: d/dx[f (x) ± g(x)] = f (x) ± g (x).

• Power rule: d/dx(αxn ) = nαxn−1 .

 
• Product rule: d/dx[f (x)g(x)] = f (x)g (x) + f (x)g(x)).
 
• Quotient rule: d/dx[f (x)/g(x)] = (g(x)f (x) − f (x)g (x))/[g(x)]2 .
 
• Chain rule: d/dx[f (g(x))] = f (g(x))g (x).

Later in this note on multivariate calculus, we are going to discuss some of the above
properties in details from a more general perspective. Until then, just remember them
so that you can use them anytime.

55
5.2 Real-Valued Functions of Several Variables
f : D → R is said to be a real-valued function if D is any set and R ⊂ R. Define
the following: x ≥ y if xi ≥ yi for every i = 1, . . . , n; and x  y if xi > yi for every
i = 1, . . . , n.

Definition 5.1 Let f : D → R, where D is a subset of Rn . Then, f is nondecreasing

if f (x) ≥ f (y) whenever x ≥ y. If, in addition, the inequality is strict whenever x  y,
then we say that f is is increasing. If, instead, f (x) > f (y) whenever x ≥ y and x
= y,
then we say that f is strongly increasing.
Rather than having a single slope, a function of n variables can be thought to have
n partial slopes, each giving only the rate at which y would change if one xi , alone, were
to change. Each of these partial slopes is called partial derivative.

Definition 5.2 Let y = f (x1 , . . . , xn ). The partial derivative of f with respect to xi is

defined as

∂f (x) f (x1 , . . . , xi + h, . . . , xn ) − f (x1 , . . . , xi , . . . , xn )

≡ lim
∂xi h→0 h

∂y/∂xi or fi (x) are used to denote partial derivatives.

5.3 Gradients
If z = F (x, y) and C is any number, we call the graph of the equation F (x, y) = C a
level curve for F . The slope of the level curve F (x, y) = C at a point (x, y) is given by
the formula
 dy ∂F (x, y)/∂x F1 (x, y)
F (x, y) = C =⇒ y = =− =−
dx ∂F (x, y)/∂y F2 (x, y)

If (x0 , y0 ) is a particular point on the level curve F (x, y) = C, the slope at (x0 , y0 ) is
−F1 (x0 , y0 )/F2 (x0 , y0 ). The equation for the tangent hyperplane T is

y − y0 = − [F1 (x0 , y0 )/F2 (x0 , y0 )] (x − x0 )

or, rearranging

F1 (x0 , y0 )(x − x0 ) + F2 (x0 , y0 )(y − y0 ) = 0.

Recalling the inner product, the equation can be written as

(F1 (x0 , y0 ), F2 (x0 , y0 )) · (x − x0 , y − y0 ) = 0

The vector (F1 (x0 , y0 ), F2 (x0 , y0 )) is said to be the gradient of F at (x0 , y0 ) is often
denoted by ∇F (x0 , y0 ) (pronounced as “nabla”). The vector (x − x0 , y − y0 ) is a vector

56
on the tangent hyperplane T which implies that ∇F (x0 , y0 ) is orthogonal to the tangent
hyperplane T at (x0 , y0 ).

Suppose more generally that F (x) = F (x1 , . . . , xn ) is a function of n variables defined

on an open set A in Rn , and let x0 = (x01 , . . . , x0n ) be a point in A. The gradient of F at
x0 is the vector


0 ∂F (x0 ) ∂F (x0 )
∇F (x ) = ,··· ,
∂x1 ∂xn

of first-order partial derivatives.

5.4 The Directional Derivative

Let z = f (x) be a function of n variables. The partial derivative ∂f /∂xi measures
the rate of change of f (x) in the direction parallel to the i-th coordinate axis. Each
partial derivative says nothing about the behavior of f in other directions. We introduce
the concept of directional derivative in order to measure the rate of change of f in an
arbitrary direction.

Consider the vector x = (x1 , . . . , xn ) and let a = (a1 , . . . , an ) ∈ Rn \{0} be a given

vector. If we move a distance ha > 0 from x in the direction given by a, we arrive at
x + ha. The average rate of change of f from x to x + ha is then (f (x + ha) − f (x))/h.
We define the derivative of f along the vector a by

 f (x + ha) − f (x)
fa (x) = lim
h→0 h
or, with components,

 f (x1 + ha1 , . . . , xn + han ) − f (x1 , . . . , xn )

fa (x1 , . . . , xn ) = lim
h→0 h
We assume that x + ha lies in the domain of f for all sufficiently small h. This is one
reason why the domain is generally assumed to be open. In particular, with ai = 1 and
aj = 0 for all j
= i, this derivative is the partial derivative of f with respect to xi .

Suppose f is C 1 in a set A 1 , and let x be an interior point in A. For an arbitrary

vector a, define the function g by

g(h) = f (x + ha) = f (x1 + ha1 , . . . , xn + han ).

1
A function f : Ên → Ê is continuously differentiable (or C 1 ) on an open set A ⊂ Ên if, for each i =
1, . . . , n, (∂f /∂xi )(x) exists for all x ∈ A and is continuous in x. f is k-times continuously differentiable
or C k on A if all the derivatives of f of order less than or equal to k(≥ 1) exist and are continuous on A.

57
 
Then, (g(h)−g(0))/h
 = (f (x+ha)−f (x))/h.
 Letting h tend to 0, we have g (0) = fa (x).
Since g (h) = ni=1 fi (x + ha)ai , g (0) = ni=1 fi (x)ai . Hence,
 

n


fa (x) = fi (x)ai = ∇f (x) · a.
i=1

This equation shows that the derivative of f along the vector a is equal to the inner

product of the gradient of f and a. If a = 1, the number fa (x) is called the directional
derivative of f at x, in the direction a.

Theorem 5.1 Suppose that f (x) = f (x1 , . . . , xn ) is C 1 in an open set A. Then, at

points x where ∇f (x) ∈ Rn \{0}, the gradient ∇f (x) = (f1 (x), . . . , fn (x)) satisfies:
1. ∇f (x) is orthogonal to the level surface through x.

2. ∇f (x) points in the direction of maximal increase of f .

3. ∇f (x) measures how fast the function increases in the direction of maximal in-
crease.
Proof of Theorem 5.1: By introducing θ as the angle between the vectors ∇f (x)
and a, we have

fa (x) = ∇f (x) · a = ∇f (x)a cos θ

Note that cos θ ≤ 1 for all θ and cos 0 = 1. So when a = 1, it follows that at points

where ∇f (x)
= 0, the number fa (x) is largest when θ = 0, i.e., when a points in the

same direction as ∇f (x), while fa (x) is smallest when θ = π, that is, cos π = −1, i.e.,
when a points in the opposite direction to ∇f (x). Moreover, it follows that the length
of ∇f (x) equals the magnitude of the maximum directional derivative. 

Theorem 5.2 (The Mean-Value Theorem) Suppose that f : Rn → R is C 1 in an

open set containing [x, y]. Then there exists a point w ∈ (x, y) such that

f (x) − f (y) = ∇f (w) · (x − y).

Proof of the mean-value theorem: We assume that the mean-value theorem for
functions of one variable is correct. Define ϕ(λ) = f (λx + (1 − λ)y). Then, using the

chain rule which we will cover later, ϕ (λ) = ∇f (λx + (1 − λ)y) · (x − y). According to
the mean-value theorem for functions of one variable, there exists a number λ0 ∈ (0, 1)

such that ϕ(1) − ϕ(0) = ϕ (λ0 ). Putting w = λ0 x + (1 − λ0 )y, the theorem follows. 

5.5 Convex Sets

Convex sets are basic building blocks in virtually every area of microeconomic the-
ory. Convexity is most often assumed to guarantee that the analysis is mathematically
tractable and that the results are clear-cut and “well-behaved.”

58
Definition 5.3 S ⊂ Rn is a convex set if for all x, y ∈ S, we have
αx + (1 − α)y ∈ S,
for all α ∈ [0, 1]
We say that z is a convex combination of x and y if z = αx + (1 − α)y for some
α ∈ [0, 1]. We have a very simple and intuitive rule defining convex sets: A set is convex
if and only if we can connect any two points in the set by a straight line that lies entirely
within the set.

Exercise 5.1 Suppose that p  0 and y ≥ 0. Let B(p, y) = {x ∈ Rn+ |p · x ≤ y} be the

budget set of the consumer. Show that B(p, y) is convex.

Theorem 5.3 Let S and T be convex sets in Rn . Then S ∩ T is a convex set.

Proof of Theorem 5.3: Let x and y be any two points in S ∩ T . Because x ∈ S ∩ T ,
we have x ∈ S and x ∈ T . Similarly, we have y ∈ S and y ∈ T . Let z = αx + (1 − α)y
for some α ∈ [0, 1] be any convex combination of x and y. z ∈ S because S is convex
and z ∈ T because T is convex. Thus, z ∈ S ∩ T . 

Exercise 5.2 Construct an example in which two sets S and T are convex but S ∪ T is
not convex.

5.5.1 Upper Contour Sets

Let u(·) : Rn → R be a utility function. Define U C(x0 ) = {x ∈ Rn+ |u(x) ≥ u(x0 )}. This
U C(x0 ) is called the upper contour set which consists of all commodity vectors x that
the individual values at least as good as x0 . In consumer theory, we usually assume that
U C(x0 ) is convex for every x0 ∈ Rn+ .

5.6 Concave and Convex Functions

A C 2 function of one variable y = f (x) is said to be concave (convex) on the interval I

if f (x) ≤ (≥) 0 for all x ∈ I.

Definition 5.4 A function f (x) = f (x1 , . . . , xn ) defined on a convex set S is concave

(convex) on S if
 
f (λx + (1 − λ)x ) ≥ (≤) λf (x) + (1 − λ)f (x )

for all x, x ∈ S and all λ ∈ [0, 1]

Definition 5.5 A function f (x) = f (x1 , . . . , xn ) defined on a convex set S is strictly

concave (convex) on S if
 
f (λx + (1 − λ)x ) > (<) λf (x) + (1 − λ)f (x )
 
for all x, x ∈ S with x
= x and all λ ∈ (0, 1)

59
5.7 Concavity/Convexity for C 2 Functions
Suppose that z = f (x) = f (x1 , . . . , xn ) is a C 2 function in an open convex set S in Rn .
The matrix

D 2 f (x) = (fij (x))n×n

is called the Hessian (matrix) of f at x, and the n determinants

 
 f11 (x) f12 (x) · · · f1r (x) 

 f21 (x) f22 (x) · · · f2r (x) 
2 
|D(r) f (x)| =  .. .. .. ..  , r = 1, . . . , n
 . . . . 
 
 fr1 (x) fr2 (x) · · · frr (x) 

are the leading principal minors of D 2 f (x) of order r. Here fij (x) = ∂ 2 f (x)/∂xi ∂xj for
any i, j = 1, . . . , r.

Theorem 5.4 (Second-Order Characterization of Concave (Convex) Functions)

Suppose that f (x) = f (x1 , . . . , xn ) is a C 2 function defined on an open, convex set S
in Rn . Let Δ2(r) f (x) denote a generic principal minor of order r in the Hessian matrix.
Then
1. f is convex in S ⇐⇒ Δ2(r) f (x) ≥ 0 for all x ∈ S and all Δ2(r) f (x), r = 1, . . . , n.

2. f is concave in S ⇐⇒ (−1)r Δ2(r) f (x) ≥ 0 for all x ∈ S and all Δ2r f (x), r =
1, . . . , n.
Proof of Theorem 5.4: (⇐=) The proof relies on the knowledge on the chain
rule (Theorem 5.15) which we are going to cover in this course. Just take it for granted
until then. Take two points x, x0 ∈ S and let t ∈ [0, 1]. Define

g(t) = f (x0 + t(x − x0 )) = f (tx + (1 − t)x0 ).

The chain rule for functions of several variables gives

n
 0 T
" 0
#
0
g (t) = (x − x ) ∇f (x + t(x − x ) = fi (x0 + t(x − x0 ))(xi − x0i )
i=1

Using the chain rule again, we get

 " #
g (t) = (x − x0 )T D2 f (x0 + t(x − x0 )) (x − x0 )
n n
= fij (x0 + t(x − x0 ))(xi − x0i )(xj − x0j )
i=1 j=1

By our hypothesis with Theorem 4.6 on quadratic forms, g (t) ≥ 0 for any t ∈ [0, 1].
This shows that g(·) is convex. In particular, we have

g(t) = g (t · 1 + (1 − t) · 0)) ≥ tg(1) + (1 − t)g(0) = tf (x) + (1 − t)f (x0 )

60
But this shows that f (·) is convex. The concavity of f easily follows by replacing f with
−f . (=⇒) Suppose f (·) is convex. According to Theorem 4.6 on quadratic forms, it
suffices to show that for all x ∈ S and all h1 , . . . , hn , we have
n
n
Q= fij (x)hi hj ≥ 0.
i=1 j=1

Now S is an open set, so if x ∈ S and h = (h1 , . . . , hn ) is an arbitrary vector, there

exists a positive number a such that x + th ∈ S for all t with |t| < a. Let I = (−a, a).
Define the function p on I by p(t) = f (x + th). Since p(·) is convex in I,
n
n

p (t) = fij (x + th)hi hj ≥ 0
i=1 j=1

for all t ∈ I. Putting t = 0, it follows that f (x) ≥ 0. This completes the proof. 

Corollary 5.1 Let z = f (x, y) be a C 2 function defined on an open convex set S ⊂ R2 .

Then,
1. f is convex ⇐⇒ f11 ≥ 0, f22 ≥ 0, and f11 f22 − (f12 )2 ≥ 0.
2. f is concave ⇐⇒ f11 ≤ 0, f22 ≤ 0, and f11 f22 − (f12 )2 ≥ 0.

Exercise 5.3 Let f (x, y) = 2x − y − x2 + 2xy − y 2 for all (x, y) ∈ R2 . Check whether f
is concave, convex, or neither.

Exercise 5.4 The CES (Constant Elasticity of Substitution) function f defined for K >
0, L > 0 by
" #−1/ρ
f (K, L) = A δK −ρ + (1 − δ)L−ρ
where A > 0, ρ
= 0, and 0 ≤ δ ≤ 1. Show that f is concave if ρ ≥ −1 and convex if
ρ ≤ −1.

Theorem 5.5 (Second-Order (Partial) Characterization of Strict Concavity)

Suppose that f (x) = f (x1 , . . . , xn ) is a C 2 function defined on an open, convex set S in
Rn . Let D(r)
2 f (x) be defined above. Then

2 f (x) > 0 for all x ∈ S and all r = 1, . . . , n =⇒ f is strictly convex.

1. D(r)

2. (−1)r D(r)
2 f (x) > 0 for all x ∈ S and all r = 1, . . . , n =⇒ f is strictly concave.

Proof of Theorem 5.5: Define the function g(·) as in the proof of Theorem 5.4
above. If the specified conditions are satisfied, the Hessian matrix D2 f (x) is positive

definite by Theorem 4.6 on quadratic forms. So, for x
= x0 , g (t) > 0 for all t ∈ [0, 1].
It follows that g(·) is strictly convex. Then, we have
g(t) = g (t · 1 + (1 − t) · 0)) > tg(1) + (1 − t)g(0) = tf (x) + (1 − t)f (x0 )
for all t ∈ (0, 1). The strict concavity of f is obtained by replacing f with −f . 

61
Corollary 5.2 Let z = f (x, y) be a C 2 function defined on an open convex set S ⊂ R2 .
Then,

1. f11 > 0 and f11 f22 − (f12 )2 > 0 =⇒ f is strictly convex.

2. f11 < 0 and f11 f22 − (f12 )2 > 0 =⇒ f is strictly concave.

Theorem 5.6 (First-Order Characterization of Concavity) Suppose that f (·) is

a C 1 function defined on an open, convex set S in Rn . Then

1. f is concave in S if and only if

n
∂f (x0 )
f (x) − f (x0 ) ≤ ∇f (x0 ) · (x − x0 ) = (xi − x0i )
∂xi
i=1

for all x, x0 ∈ S.

2. f (·) is strictly concave iff the above inequality is always strict when x

= x0 .

3. The corresponding result for convex (strictly convex) functions is obtained by chang-
ing ≤ to ≥ (< to >) in the above inequality.

Proof of Theorem 5.6: (1) (=⇒) Let x, x0 ∈ S. Since f is concave,

λf (x) + (1 − λ)f (x0 ) ≤ f (λx + (1 − λ)x0 )

for all λ ∈ (0, 1). Rearranging the above inequality, for all λ ∈ (0, 1), we obtain

f (x0 + λ(x − x0 )) − f (x0 )

f (x) − f (x0 ) ≤ (∗)
λ
Let λ → 0. The right hand side of (∗) then approaches ∇f (x0 ) · (x − x0 ). (⇐=) Let
x, x0 ∈ S and λ ∈ (0, 1). Define z = λx + (1 − λ)x0 . Notice that z ∈ S because S is
convex. By our hypothesis, we have

f (x) − f (z) ≤ ∇f (z) · (x − z) (i)

f (x0 ) − f (z) ≤ ∇f (z) · (x0 − z) (ii)

Multiplying the inequality in (i) by λ > 0 and the inequality in (ii) by 1 − λ > 0, we
obtain
! " #
λ (f (x) − f (z)) + (1 − λ) f (x0 ) − f (z) ≤ ∇f (z) · λ(x − z) + (1 − λ)(x0 − z) (iii)

Here λ(x − z) + (1 − λ)(x0 − z) = λx + (1 − λ)x0 − z = 0, so the right hand side of (iii)

is 0. Thus, rearranging (iii) gives

λf (x) + (1 − λ)f (x0 ) ≤ f (z) = f (λx + (1 − λ)x0 )

62
because z = λx + (1 − λ)x0 . This shows that f is concave. (2) (=⇒) Suppose that
f is strictly concave in S. Then, inequality (∗) is strict for x
= x0 . (⇐=) With z =
x0 + λ(x − x0 ), we have

f (z) − f (x0 ) ∇f (x0 ) · (z − x0 )

f (x) − f (x0 ) < ≤ = ∇f (x0 ) · (x − x0 ).
λ λ
where we used the inequality in (1), which we have already proved, and the fact that
z − x0 = λ(x − x0 ). This shows that the inequality in (1) holds with strict inequality.
(3) This part is trivial. Do you agree with me? 

5.7.1 Jensen’s Inequality

Theorem 5.7 (Jensen’s Inequality) A function f (·) is concave on a convex set S in
Rn if and only if

f (λ1 x1 + · · · λn xn ) ≥ λ1 f (x1 ) + · · · λn f (xn )

n
holds for all x1 , . . . , xn ∈ S, and for all λi ≥ 0 for all i = 1, . . . , n with i=1 λi = 1.
Proof of Jensen’s Inequality: For any k ≥ 2, we propose the hypothesis H(k) as
follows.
 k  k

H(k) : f λk xh ≥ λh f (xh )
h=1 h=1

H(2) is true because it is indeed the definition of concavity of f . Now, we will show that
H(k) =⇒ H(k + 1). We execute a series of computations below.
k+1   k $ k % 
λh
f λh xh = f λh k xh + λk+1 xk+1
h=1 h=1 h=1 h=1 λh
 k   k 
λh
≥ λh f k xh + λk+1 f (xk+1 ) (because of H(2))
h=1 h=1 h=1 λh
 k  k
λh
≥ λh · k f (xh ) + λk+1 f (xk+1 )
h=1 h=1 h=1 λh
(because H(k) is true under the inductive hypothesis.)
k+1

= λh f (xh ). 
h=1

One can extend Jensen’s inequality to the continuum. Let X be a random variable
which takes values on the real line R. Define g : R → R to be a probability density
function. Then, continuous version of Jensen’s inequality is given:

63
 Jensen’s Inequality (Continuum Version): A function f (·) is concave on R if
and only if

& ∞  & ∞
f f (x)g(x)dx ≥ f (x)g(x)dx
−∞ −∞

for any probability density function g(·).

5.8 Quasiconcave and Quasiconvex Functions

Definition 5.6 A function f , defined over a convex set S ⊂ Rn , is quasiconcave if
the upper level set Pα = {x ∈ S|f (x) ≥ α} is convex for each α ∈ R. We say that
f is quasiconvex if −f is quasiconcave. So, f is quasiconvex iff the lower level set
P α = {x ∈ S|f (x) ≤ α} is convex for each α ∈ R.

Proposition 5.1 If f (·) is concave, then it is quasiconcave. Similarly, if f (·) is convex,

then it is quasiconvex.

Exercise 5.5 Prove Proposition 5.1.

Theorem 5.8 Let f (·) be a function of n variables defined on a convex set S in Rn .

Then, f is quasiconcave if and only if either of the following conditions is satisfied for

all x, x ∈ S and all λ ∈ [0, 1],
 
1. f (λx + (1 − λ)x ) ≥ min{f (x), f (x )}
 
2. f (x ) ≥ f (x) =⇒ f (λx + (1 − λ)x ) ≥ f (x)

Proof of Theorem 5.8: (1) (=⇒) Suppose that f (·) is quasiconcave. Let x, x ∈ S

and λ ∈ [0, 1], and define a = min{f (x), f (x )}. Then,

x, x ∈ Pa = {x ∈ S|f (x) ≥ a}

Since Pa is convex by our hypothesis, λx + (1 − λ)x ∈ Pa for any λ ∈ [0, 1]. This implies
 
that f (λx + (1 − λ)x ) ≥ a = min{f (x), f (x )}. (⇐=) Suppose that the inequality in (1)
is valid and let a be an arbitrary number. We must show that Pa is convex. Take any
 
arbitrary points x, x ∈ Pa . Then, f (x) ≥ a and f (x ) ≥ a. Also, for all λ ∈ (0, 1), the
inequality in (1) implies that
' 
( 
f λx + (1 − λ)x ≥ min{f (x), f (x )}


Thus, λx + (1 − λ)x ∈ Pa . This proves that Pa is convex. We leave the rest of the proof
as an exercise. 

Exercise 5.6 Prove the second part of Theorem 5.8.

A function : R → R is said to be strictly increasing if F (x) > F (y) whenever x > y.

64
Theorem 5.9 (Quasiconcavity is preserved under positive monotone transformation)
Let f (·) be defined on a convex set S in Rn and let F be a function of one variable whose
domain includes f (S). If f (·) is quasiconcave (quasiconvex) and F is strictly increasing,
then F (f (·)) is quasiconcave (quasiconvex).
Proof of Theorem 5.9: Suppose f (·) is quasiconcave. Using the previous theorem
(Theorem 5.8), we must have
' 
( 
f λx + (1 − λ)x ≥ min{f (x), f (x )}.

Since F (·) is strictly increasing,

' ' 
(( ' 
( 
F f λx + (1 − λ)x ≥ F min{f (x), f (x )} = min{F (f (x)), F (f (x ))}.

It follows that F ◦ f is quasiconcave. The argument in the quasiconvex case is entirely

similar, replacing ≥ with ≤ and min with max. 

Definition 5.7 A function f (·) defined on a convex set S ⊂ Rn is said to be strictly

quasiconcave if
 
f (λx + (1 − λ)x ) > min{f (x), f (x )}
 
for all x, x ∈ S with x
= x and all λ ∈ (0, 1). The function f is strictly quasiconvex
if −f is strictly quasiconcave.
n
Exercise 5.7 The Cobb-Douglas function f (x) = Axα1 1 · · · xαn , defined for x1 > 0, . . . , xn >
0, with A > 0 and αi > 0 for all i = 1, . . . , n. Show the following.
1. f (·) is quasiconcave for all α1 , . . . , αn .

2. f (·) is concave for α1 + · · · + αn ≤ 1.

3. f (·) is strictly concave for α1 + · · · + αn < 1.

Theorem 5.10 (First-Order Characterization of Quasiconcavity) Let f (·) be a

C 1 function of n variables defined on an open convex set S in Rn . Then f (·) is quasi-
concave on S if and only if for all x, x0 ∈ S,
n

0 0 0 ∂f (x0 )
f (x) ≥ f (x ) =⇒ ∇f (x ) · (x − x ) = (xi − x0i ) ≥ 0.
∂xi
i=1

Proof of Theorem 5.10: (=⇒) Suppose f (·) is quasiconcave. Let x, x0 ∈ S and

define the function g(·) on [0, 1] by
! !
g(t) = f (1 − t)x0 + tx = f x0 + t(x − x0 ) .

Then, using the chain rule (Theorem 5.15) which will be shown later, we have

g (t) = ∇f (x0 + t(x − x0 )) · (x − x0 ).

65
Suppose f (x) ≥ f (x0 ). By Theorem 5.8, g(t) ≥ g(0) for all t ∈ [0, 1]. For any t ∈ (0, 1],
we have
g(t) − g(0)
≥ 0.
t
Letting t → 0, we obtain
g(t) − g(0) 
lim = g (0) ≥ 0.
t→0 t
This implies

g (0) = ∇f (x0 ) · (x − x0 ) ≥ 0

(⇐=) We will be satisfied with the figure for this part. 

The content of Theorem 5.10 is that for any quasiconcave function f (·) and any pair
of points x and x0 with f (x) ≥ f (x0 ), the gradient vector ∇f (x0 ) and the vector (x − x0 )
must form an acute angle.

Theorem 5.11 (Second-Order Characterization of Quasiconcavity) Let f (·) be

a C 2 function defined on an open, convex set S in Rn . Then, f (·) is quasiconcave if
and only if, for every x ∈ S, the Hessian matrix D 2 f (x) is negative semidefinite in the
subspace {z ∈ Rn |∇f (x) · z = 0}, that is,

z T D2 f (x)z ≤ 0 whenever ∇f (x) · z = 0

for every x ∈ S. If the Hessian matrix D2 f (x) is negative definite in the subspace
{z ∈ Rn |∇f (x) · z = 0} for every x ∈ S, then f (·) is strictly quasiconcave.
Proof of Theorem 5.11: (=⇒) Suppose f (·) is quasiconcave. Let x ∈ S. Choose
  
x ∈ S such that ∇f (x) · (x − x) = 0. Since f (·) is quasiconcave, f (x ) ≤ f (x). To see
this, draw the figure. Then,
 
f (x ) − f (x) ≤ ∇f (x) · (x − x) = 0

By Theorem 5.6, f (·) is concave in the subspace for which ∇f (x) · (x − x) = 0. With
Theorem 5.4 and Theorem 4.6, concavity of f (·) is equivalent to negative semidefiniteness
of the Hessian matrix. Then, the conclusion follows. (⇐=) This proof is based on “A
Characterization of Quasi-Concave Functions,” by Kiyoshi Otani in Journal of Economic
 
Theory, vol 31, (1983), 194-196. Let x, x ∈ S such that f (x ) ≥ f (x). This choice entails
no loss of generality. For λ ∈ [0, 1], define

g(λ) = f (x + λ(x − x)).
 
Note that g(0) = f (x), g(1) = f (x ), and g(1) ≥ g(0) because f (x ) ≥ f (x) by our
hypothesis. What we want to show is that g(λ) ≥ g(0) for any λ ∈ [0, 1]. By the

mean-value theorem (Theorem 5.2), there exists λ0 ∈ (0, 1) such that g (λ0 ) = 0. Let

66
   
x0 = λ0 x+ (1− λ0 )x . Then, g (λ0 ) = 0 ⇔ ∇f (x0 )·(x − x) = 0. Assume that ∇f (x)
= 0
for any x ∈ S. This strikes us as being innocuous. Let p denote ∇f (x0 ) for notational
simplicity.

By pT · p > 0, there exists a C 2 function β : R → R for sufficiently small |α| > 0 such
that
β(0) = 0 and
' 
(
f β(α)p + α(x − x) + x0 = f (x0 ) for any small α

Again, for notional simplicity, we denote β(α)p + α(x − x) + x0 by z(α). By differ-
entiating f (z(α)) = f (x0 ), we have
)  
*
∇f (z(α)) β (α)p + (x − x) = 0 (∗)

and by further differentiating, we have

)  
*T )  
* 
β (α)p + (x − x) D2 f (z(α)) β (α)p + (x − x) + ∇f (z(α))β (α)p = 0. (∗∗)
 
Since pT · (x − x) = 0 and ∇f (z(0))p ∈ Rn \{0}, (∗) implies β (0) = 0. Moreover, our
hypothesis requires
)  
*T )  
*
β (α)p + (x − x) D2 f (z(α)) β (α)p + (x − x) ≤ 0.

Then, we must have β (α)∇f (z(α))p ≥ 0. When α is sufficiently close to zero, z(α) is

very close to x0 and so ∇f (z(α))p > 0 because p ∈ Rn \{0}. Then, β (α) ≥ 0 for α with
|α| sufficiently small.

For sufficiently small |λ − λ0 |, we have

' 
(
∇f (t − t0 )(x − x) + x0 p > 0

because ∇f (x0 )T p > 0 and (t−t0 )(x −x)+x0 is very close to x0 . Hence, for λ sufficiently
close to λ0 , we have

' 
(
g(λ) = f (x + λ(x − x)) = f (λ − λ0 )(x − x) + x0
' 
(
≤ f β(λ − λ0 )p + (λ − λ0 )(x − x0 ) + x0 = f (x0 ) = g(λ0 )

because β(λ − λ0 ) ≥ 0 for λ sufficiently close to λ0 . Accordingly, g(λ0 ) does not have an
interior minimum in [0, 1], unless it is constant. Hence, g(λ) ≥ g(0) for any λ ∈ [0, 1].
The last step is based on Corollary 4.3 in “Nine Kinds of Quasiconcavity and Concavity,”
by Diewert, Avriel, and Zang in Journal of Economic Theory, vol 25, (1981), 397-420.

Corollary 4.3 (Diewert, Avriel, and Zang (1981)): A differentiable function f

defined over an open S is quasiconcave if and only if, for any x0 ∈ S and any v ∈ R with
v T v = 1,

67
v T ∇f (x0 ) = 0 implies g(t) ≡ f (x0 + tv) does not attain a (semistrict) local minimum at
t = 0.

This completes the proof. 

Theorem 5.12 (A Characterization through Bordered Hessian) Let f be a C 2

function defined in an open, convex set S in Rn . Define the bordered Hessian determi-
nants Br (x) as follows: for each r = 1, . . . , n,
 
 0 f1 (x) · · · fr (x) 

 f1 (x) f11 (x) · · · f1r (x) 
 
Br (x) =  . .. .. .. .
 .. . . . 
 
 fr (x) fr1 (x) · · · frr (x) 

Then,

1. A necessary condition for f to be quasiconcave is that (−1)r Br (x) ≥ 0 for all x ∈ S

and all r = 1, . . . , n.

2. A sufficient condition for f to be strictly quasiconcave is that (−1)r Br (x) > 0 for
all x ∈ S and all r = 1, . . . , n.

5.9 Total Differentiation

Consider functions that map points (vectors) in Rn to points (vectors) in Rm . Such
functions are often called transformations (operators). A transformation f : Rn → Rm
is said to be linear if

f (x1 + x2 ) = f (x1 ) + f (x2 ) and f (αx1 ) = αf (x1 )

for all x1 , x2 ∈ Rn and all scalars α ∈ R. Our knowledge on linear algebra tells us that
for every linear transformation f : Rn → Rm , there is a unique m × n matrix A such
that f (x) = Ax for all x ∈ Rn .
⎛ (1) ⎞ ⎛ ⎞⎛ ⎞
f (x) a11 a12 · · · a1n x1
⎜ .. ⎟ ⎜ .. .. .. .. ⎟ ⎜ x2 ⎟
⎜ . ⎟ ⎜ . . . . ⎟ ⎜ ⎟
⎜ ⎟ ⎜ ⎟⎜ ⎟
⎜ f (j)(x) ⎟ = ⎜ aj1 aj2 · · · ajn ⎟ ⎜ ... ⎟
⎜ ⎟ ⎜ ⎟⎜ ⎟
⎜ .. ⎟ ⎜ .. .. .. .. ⎟ ⎜ .. ⎟
⎝ . ⎠ ⎝ . . . . ⎠ ⎝ . ⎠
(m)
f (x) am1 am2 · · · amn xn

In particular,
n

f (j)(x) = aj1 x1 + aj2 x2 + · · · ajn xn = aji xi .
i=1

68
5.9.1 Linear Approximations and Differentiability
If a one variable function f is differentiable at a point x0 , then the linear approximation
to f around x0 is given by

f (x0 + h) ≈ f (x0 ) + f (x0 )h,

for small values of h. Here ≈ stands for “approximation.” This is useful because the
approximation error is defined by

O(h) = true value - approximate value = f (x0 + h) − f (x0 ) − f (x0 )h

becomes negligible for sufficiently small h. Namely O(h) → 0 as h → 0. More impor-

tantly, however, O(h) also becomes small in comparison with h - that is,


O(h) f (x0 + h) − f (x0 )  0
lim = lim − f (x ) = 0.
h→0 h h→0 h

Moreover, f (·) is differentiable at x0 if and only if there exists a number c ∈ R such that

f (x0 + h) − f (x0 ) − ch
lim = 0.
h→0 h

If such a number c ∈ R exists, it is unique and c = f (x0 ). These arguments can be gen-
eralized straightforwardly to higher dimensional spaces. In particular, a transformation
f (·) is differentiable at a point x0 if it admits a linear transformation around x0 :

Definition 5.8 If f : A → Rm is a transformation defined on a subset A of Rn and

x0 is an interior point of A, then f is said to be differentiable at x0 if there exists an
m × n matrix C such that

f (x0 + h) − f (x0 ) − Ch

lim =0
h→0∈Ên h

If such a matrix C exists, it is called the total derivative of f (·) at x0 , and is denoted
by Df (x0 ).
If I restrict attention to real-valued functions, the following theorem establishes an
equivalence between directional derivative along every vector and total differentiation.

Theorem 5.13 If f : A → R is defined on a subset A of Rn and f is differentiable


at an interior point x ∈ A, then f has a derivative fa (x) along every n-vector a, and

fa (x) = ∇f (x) · a.
Proof of Theorem 5.12: The derivative along a is


 f (x + ha) − f (x) − ∇f (x) · ah
fa (x) = lim + ∇f (x) · a = 0 + ∇f (x) · a
h→0 h

69
 j


In particular, if ej = (0, . . . , 1 , . . . , 0) is the jth standard unit vector in Rn , then
∇f (x) · ej is the partial derivative ∂f (x)/∂xj = fj (x) with respect to the jth variable.
On the other hand, ∇f (x) · ej is the jth component of ∇f (x). Hence, ∇f (x) is the row
vector
∇f (x) = (∇f (x) · e1 , . . . , ∇f (x) · en ) = (f1 (x), . . . , fn (x)) .
Suppose that I am interested in checking if a given transformation f : Rn → Rm is
differentiable. Then, the next theorem shows that it suffices to check if each component
real-valued function f j : Rn → R is differentiable (j = 1, . . . , m). 2

Theorem 5.14 A transformation f = (f 1 , . . . , f m ) from a subset A of Rn into Rm is

differentiable at an interior point x ∈ A if and only if each component function f j : A →
R, j = 1, . . . , m, is differentiable at x. Moreover,
⎛ ⎞ ⎛ ∂f 1 ∂f 1 ∂f 1
⎞
∇f (1) (x) ∂x1 (x) ∂x2 (x) · · · ∂xn (x)
⎜ ∇f (2) (x) ⎟ ⎜ ⎜ ∂f 2 ∂f 2
· · · ∂f 2 ⎟
⎟
⎜ ⎟ ⎜ ∂x1 (x) ∂x2 (x) ∂xn (x) ⎟
Df (x) = ⎜ .. ⎟=⎜ . . . . ⎟
⎝ . ⎠ ⎝ .. .. .. .. ⎠
∇f (m) (x) ∂f m ∂f m ∂f m
∂x1 (x) ∂x2 (x) · · · ∂xn (x)

This is called the Jacobian matrix of f (·) at x. Its rows are the gradients of the
component functions of f (·).
Proof of Theorem 5.13: Let C be an m × n matrix and let O(h) = f (x + h) −
f (x) − Ch where h ∈ Rn .
⎛ ⎞ ⎛ (1) ⎞ ⎛ ⎞⎛ ⎞
O1 (h) f (x + h) − f (1) (x) c11 c12 · · · c1n h1
⎜ O2 (h) ⎟ ⎜ f (2) (x + h) − f (2) (x) ⎟ ⎜ c21 c22 · · · c2n ⎟ ⎜ h2 ⎟
⎜ ⎟ ⎜ ⎟ ⎜ ⎟⎜ ⎟
⎜ .. ⎟=⎜ .. ⎟ − ⎜ .. .. .. .. ⎟ ⎜ .. ⎟ .
⎝ . ⎠ ⎝ . ⎠ ⎝ . . . . ⎠ ⎝ . ⎠
Om (h) (m) (m)
f (x + h) − f (x) cm1 cm2 · · · cmn hn

The j-th component of O(h), j = 1, . . . , m, is Oj (h) − f (j) (x + h) − f (j) (x) − Cj h, where

Cj is the j-th row of C. For each j,
|Oj (h)| ≤ O(h) ≤ |O1 (h)| + · · · |Om (h)|
It follows that
O(h) |Oi (h)|
lim = 0 ⇐⇒ lim = 0 for all i = 1, . . . , m
h→0 h h→0 h

Hence, f (·) is differentiable at x if and only if each f (j) is differentiable at x. Also, the
j-th row of the matrix C = Df (x) is the derivative of f (j), that is, Cj = ∇f (j) (x). 

The next theorem confirms our intuition about differentiation: If a transformation

(or function) is differentiable, it is more than continuous.
2
See the similar argument in Theorem 3.14 in which we show that a function f : Ên → Êm is
continuous if and only if each component function f j : Ên → Ê is continuous (j = 1, . . . , m).

70
Theorem 5.15 (Differentiability =⇒ Continuity) If a transformation f from A ⊂
Rn into Rm is differentiable at an interior point x0 ∈ A, then f is continuous.
Proof of Theorem 5.14: Let C = Df (x). Then, for small but nonzero h ∈ Rn , the
triangle inequality yields

f (x0 + h) − f (x0 ) = f (x0 + h) − f (x0 ) + Ch − Ch

≤ f (x0 + h) − f (x0 ) + Ch + Ch (∵ Minkowski inequality)


f (x0 + h) − f (x0 ) + Ch
= h + Ch
h

Since f is differentiable at x0 ,

f (x0 + h) − f (x0 ) + Ch

→ 0 as h → 0
h
Ch → 0 as h → 0

Hence, f (x0 + h) → f (x0 ) as h → 0. 

The next theorem shows that the order of two operations do not matter for the final
product: (1) Construct a composite mapping and differentiate it; and (2) Differentiate
each mapping and constructs a composite of two derivatives.

Theorem 5.16 (The Chain Rule) Suppose f : A → Rm and g : B → Rp are defined

on A ⊂ Rn and B ⊂ Rm , with f (A) ⊂ B, and suppose that f and g are differentiable at
x and f (x), respectively. Then, the composite transformation g ◦ f : A → Rp defined by
(g ◦ f )(x) = g(f (x)) is differentiable at x, and

D(g ◦ f )(x) = Dg(f (x)) × Df (x)

        
p×n p×m m×n

Proof of the Chain Rule: Define

ef (h)
k(h) = f (x + h) − f (x) = Df (x)h + ef (h), where → 0 as h → 0
h

Also,

eg (k)
g(f (x + h)) − g(f (x)) = Dg(f (x))k + eg (k), where → 0 as k → 0
k

Note that there exits some fixed constant K such that k(h) ≤ Kh for all small h.
Otherwise, f and g are not differentiable. Note also that for all ε > 0, eg (k) < εk
for k small because g is differentiable. Thus, when h is small, we can summarize

eg (k(h)) < εk(h) ≤ εKh

71
Hence,
eg (k(h))
→ 0 as h → 0
h
Then, we execute a series of computation below.

e(h) = g (f (x) + k(h)) − g(f (x)) − Dg(f (x))Df (x)h

= D(g ◦ f )(x)k(h) + eg (k(h)) − Dg(f (x))Df (x)h
= D(g ◦ f )(x) (k(h) − Df (x)h) + eg (k(h))
= D(g ◦ f )(x)ef (h) + eg (k(h)) (∵ h(k) = Df (x)h + ef (h))

And, moreover,
e(h) 1
= D(g ◦ f )(x)ef (h) + eg (k(h))
h h
D(g ◦ f )(x)ef (h) eg (k(h))
≤ + (∵ triangle inequality)
h h
ef (h) eg (k(h))
≤ D(g ◦ f )(x) + (∵ Cauchy-Schwartz inequality)
h h
Since ef (h)/h → 0 and eg (k(h))/h → 0 as h → 0, we conclude that e(h)/h →
0 as h → 0. 

5.10 The Inverse of a Transformation

Consider a transformation f : A → B where A ⊂ Rn and B ⊂ Rm . Suppose the range of
 
f is the whole of B, i.e., f (A) = B. Recall that f is one-to-one if f (x) = f (x ) =⇒ x = x .
In sum, f is bijective. In this case, for each point y ∈ B there is exactly one point x ∈ A
such that f (x) = y, and the inverse of f is the transformation f −1 : B → A which maps
each y ∈ B to precisely that point x ∈ A for which f (x) = y.

If f : U → V and g : V → U are differentiable and mutually inverse transformations

between open sets U and V in Rn , then g ◦ f is the identity transformation on U , and
therefore D(g ◦ f )(x) = In for all x ∈ U . The chain rule then gives Dg(f (x))Df (x) = In .
This means that the Jacobian matrix Df (x) must be nonsingular, so |Df (x)|
= 0. Also,
Dg(f (x)) is the inverse matrix Df (x)−1 .

Theorem 5.17 (Inverse Function Theorem) Consider a transformation f = (f1 , . . . , fn )

from A ⊂ Rn into Rn and assume that f is C k (k ≥ 1) in an open set containing
x0 = (x01 , . . . , x0n ). Furthermore, suppose that |Df (x)|
= 0 for x = x0 . Let y 0 = f (x0 ).
Then, there exists an open set U around x0 such that f maps U one-to-one onto an open
set V around y 0 , and there is an inverse mapping g = f −1 : V → U which is also C k .
Moreover, for all y ∈ V , we have

Dg(y) = Df (x)−1 , where x = g(y) ∈ U.

72
 “Sketch” of the Proof of Inverse Function Theorem: The proof consists of
5 steps. However, the proof of each step will be either briefly sketched or completely
skipped due to its technical difficulty. For simplicity, we assume that x0 = 0 and f (x0 ) =
y0 = 0.

Step 1: There is no loss of generality to assume that Df (x) = In

Proof of Step 1: Let Df (0) = A. Since A is non-singular, A−1 exists. Let g : Rn →

Rn be a linear mapping associated with A−1 . That is, g(x) = A−1 x for any x ∈ Rm .
Note that g(0) = 0 and Dg(0) = A−1 .
The Jacobian matrix of f ◦ g : Rn → Rn is given as D(f ◦ g)(0) = Df (0)Dg(0) =
AA = In . If we can show that f ◦ g is C r , so is f . This is because g is a linear map
−1

associated with A−1 and so is C k . Therefore, we can rather talk about f ◦ g instead of
f so that Df (0) = In can be assumed with no loss of generality. 

Step 2: There exists an open set U containing 0 such that f |U : U → Rn is one to

one.

Proof of Step 2:
Step 3: There exists an open set V such that f |U : U → V is onto. That is, for any
0 ∈ V , there exists 0 ∈ U such that f (0) = 0.

Step 4: f −1 |V : V → U is one-to-one and onto.

Step 5: f −1 |V : V → U is C k . 

5.11 Implicit Function Theorems

f1 (x1 , x2 , . . . , xn , y1 , y2 , . . . , ym ) = 0
· · · · · · · · · · · · · · · · · · · · · · · · · · · · · · · · · · · · ⇐⇒ f (x, y) = 0
fm (x1 , x2 , . . . , xn , y1 , y2 , . . . , ym ) = 0

with f = (f1 , . . . , fm ) , x = (x1 , . . . , xn ), and y = (y1 , . . . , ym ).

⎛ ⎞
∂f1 /∂x1 · · · ∂f1 /∂xn
⎜ .. .. .. ⎟
Dx f (x, y) = ⎝ . . . ⎠.
∂fm /∂x1 · · · ∂fm /∂xn

Theorem 5.18 (Implicit Function Theorem) Suppose f = (f 1 , . . . , f m ) is C 1 in

an open set A ⊂ Rn+m , and consider the vector equation f (x, y) = 0, where x ∈ Rn
and y ∈ Rm . Let (x0 , y 0 ) be an interior point of A satisfying f (x, y) = 0. Suppose that

73
the Jacobian determinant of f with respect to y is different from 0 at (x0 , y 0 ) – i.e.,
|Dy f (x, y)|
= 0 at (x, y) = (x0 , y 0 ). Then, there exist open balls B1 and B2 around x0
and y 0 , respectively, with B1 × B2 ⊂ A, such that |Dy f (x, y)|
= 0 in B1 × B2 , and such
that for each x ∈ B1 , there is a unique y ∈ B2 with f (x, y) = 0. In this way, y is defined
on B1 as a C 1 function g(·) of x. The Jacobian matrix Dg(x) can be found by implicit
differentiation of f (x, y) = 0, and
⎛ ⎞
∂y1 /∂x1 ∂y1 /∂x2 · · · ∂y1 /∂xn
dy ⎜ ⎜ ∂y2 /∂x1 ∂y2 /∂x2 · · · ∂y2 /∂xn ⎟
⎟
=⎜ .. .. . . ⎟ = Dg(x) = −[Dy f (x, y)]−1 Dx f (x, y)
dx ⎝ . . .. .. ⎠
∂ym /∂x1 ∂ym /∂x2 · · · ∂ym /∂xn

Proof of Implicit Function Theorem: We define the norm of vectors in Rn as

follows.

x ≡ max |xi |

1≤i≤n

This is the norm we discuss for the implicit function theorem. The proof relies on the
following three lemmas (Lemmas 5.1, 5.2, and 5.3). We will not provide their proofs
here.

Lemma 5.1 Let K be a compact set in Rn . Let {hk (x)}Kk=1 be a sequence of continuous
m
functions K → R . Suppose that for any ε > 0, there exists a number N ∈ N such that

max hm (x) − hn (x) < ε

x∈K

for all m, n > N . Then, there exists a unique continuous function h : K → Rm such that
 +
lim max hk (x) − h(x) = 0
k→∞ x∈K

With Weierstrass’s Theorem (Theorem 3.19) and the concept of Cauchy sequence in
Rn (Definition 3.6), the above lemma should be easy to be established. For the next
lemma, define the following.
 
Dα = x ∈ Rn x − x0  ≤ α
 
Dβ = y ∈ Rm y − y 0  ≤ β

Let ξ(x, y) be a continuous mapping from Dα × Dβ to R with the property that

ξ(x0 , y 0 )
= 0. Notice that Dα × Dβ is a compact set by construction.

Lemma 5.2 (Lipschitz Continuity) There exists a number K ∈ (0, 1) such that, for

all y, y ∈ Dβ ,
 
ξ(x, y) − ξ(x, y ) < Ky − y .

74
 We argue that Lemma 5.2 enables us to construct a sequence of continuous functions
needed for Lemma 5.1. Again, we take Lemma 5.2 for granted. Let y0 (x) = y 0 . Define
yk+1 (x) = y0 + ξ(x, yk (x)) for k ≥ 0. Since ξ(x0 , y 0 ) = 0, we have

|ξ(x, y 0 )| = (1 − K)β

for x ∈ Dα with α > 0 sufficiently small. We execute a series of computations below.

|yk+1 (x) − y 0 | = |ξ(x, yk (x))|

= |ξ(x, yk (x)) − ξ(x, y 0 ) + ξ(x, y 0 )|
≤ |ξ(x, yk (x)) − ξ(x, y 0 )| + |ξ(x, y 0 )|
< K|yk (x) − y 0 | + (1 − K)β

Fix m ∈ N.

yk+m (x) − yk (x) = ξ(x, yk+m−1 (x)) − ξ(x, yk−1 (x))

≤ Kyk+m−1 (x) − yk−1 (x)
≤ K k+1 ym (x) − y 0  → 0 as k → ∞

This means that maxx∈Dα ym (x) − yn (x) → 0 as m, n → ∞.

Lemma 5.3 Let ξ(x, y) be a continuous mapping from Dα × Dβ to Rm with the property
that ξ(x0 , y 0 ) = 0. Furthermore, There exists a number K ∈ (0, 1) such that, for all

y, y ∈ Dβ ,
 
ξ(x, y) − ξ(x, y ) < Ky − y .

Then, there exists a unique continuous mapping ϕ : Dα → Rm for which

ϕ(x) − y 0 = ξ(x, ϕ(x))

for x ∈ Dα with α > 0 sufficiently small.

With the help of all three lemmas above, we will complete the proof of Implicit
Function Theorem.

Define a function g(x, y) : Dα × Dβ → Rm satisfying the following equation.

m×1
  
f (x, y) = Dy f (x0 , y 0 ) (y − y 0 ) + g(x, y)
        
m×1 m×m m×1

Since |Dy f (x0 , y 0 )|

= 0, we have
" #−1
y − y 0 = − Dy f (x0 , y 0 ) g(x, y)

75
By construction of g, we have g(x0 , y 0 ) = 0 and |Dy g(x0 , y 0 )|
= 0. Define
" #−1
ξ(x, y) = − Dy f (x0 , y 0 ) g(x, y)

Note also that ξ(x0 , y 0 ) = 0. By the mean-value theorem (Theorem 5.2),

 " #−1 
ξ(x, y) − ξ(x, y ) = − Dy f (x0 , y 0 ) [Dy g(x, y)](y − y )
" #−1 " # 
= − Dy f (x0 , y 0 ) Dy f (x, y) − Dy f (x0 , y 0 ) (y − y )

= Im − [Dy f (x0 , y 0 )]−1 Dy f (x, y) (y − y 0 )

If we choose α > 0 small enough so that x is very close to x0 , i.e., x ∈ Dα , there exists
K ∈ (0, 1) such that

ξ(x, y) − ξ(x, y ) < KIm (y − y 0 )
= Ky − y 0 

Now, we can take two open sets B1 and B2 small enough so that B1 ⊂ Dα and B2 ⊂ Dβ
needed for the theorem. Then, we can use Lemma 5.3 which completes the proof. 

Corollary 5.3 (A Version of Implicit Function Theorem) Suppose f : R2 → R is

C 1 in an open set A containing (x0 , y0 ), with f (x0 , y0 ) = 0 and ∂f (x0 , y0 )/∂y
= 0. Then,
there exists an interval Ix = (x0 − δ, x0 + δ) and an interval Iy = (y0 − ε, y0 + ε) (with
δ > 0 and ε > 0) such that Ix × Iy ⊂ A and:

1. for every x ∈ Ix , the equation f (x, y) = 0 has a unique solution in Iy which defines
y as a function y = ϕ(x) in Iy ;

2. ϕ is C 1 in Ix = (x0 − δ, x0 + δ), with derivative

dy  ∂f (x, ϕ(x))/∂x
= ϕ (x) = −
dx ∂f (x, ϕ(x))/∂y

Exercise 5.8 The point P = (x, y, z, u, v, w) = (1, 1, 0, −1, 0, 1) satisfies all the equa-
tions

y 2 − z + u − v − w3 = −1
−2x + y − z 2 + u + v 3 − w = −3
x2 + z − u − v + w 3 = 3

Using the implicit function theorem, find du/dx, dv/dx, and dw/dx at P .

76
Chapter 6

Static Optimization

6.1 Unconstrained Optimization

6.1.1 Extreme Points
Let f (·) be a real-valued function of n variables x1 , . . . , xn defined on a set S in Rn .
Suppose that the point x∗ = (x∗1 , . . . , x∗n ) belongs to S and that the value of f at x∗ is
greater than or equal to the values attained by f at all other points x = (x1 , . . . , xn ) ∈ S.
Thus,

f (x∗ ) ≥ f (x) for all x ∈ S (∗)

Here x∗ is called a (global) maximal point for f in S and f (x∗ ) is called the maximum
value. If the inequality (∗) is strict for all x
= x∗ , then x∗ is a strict maximum point
for f (·) in S. We define (strict) minimum point and minimum value by reversing the
inequality sign in (∗). As collective names, we use extreme points and extreme values to
indicate both maxima or minima.

Theorem 6.1 Let f (·) be defined on a set S in Rn and let x∗ = (x∗1 , . . . , x∗n ) be an
interior point in S at which f (·) has partial derivatives. A necessary condition for x∗ to
be an extreme point for f is that x∗ is a stationary point for f (·) – that is, it satisfies
the equations

∂f (x)
∇f (x) = 0 ⇐⇒ = 0, for i = 1, . . . , n
∂xi
Proof of Theorem 6.1: Suppose, on the contrary, that x∗ is a maximum point but
not a stationary point for f (·). Then, there is no loss of generality to assume that there
exists at least i such that fi (x) > 0. Define x∗∗ = (x∗1 , . . . , x∗i + ε, . . . , x∗n ). Since x∗ is an
interior point in S, one can make sure that x∗∗ ∈ S by choosing ε > 0 sufficiently small.
Then,

f (x∗∗ ) ≈ f (x∗ ) + ∇f (x) · (0, . . . , 0, 

ε , 0, . . . , 0) > f (x∗ ).
i

77
However, this contradicts the hypothesis that x∗ is a maximum point for f (·). 

The next theorem clarifies under what conditions, the converse of the previous theo-
rem (Theorem 6.1) is established.

Theorem 6.2 Suppose that the function f (·) is defined in a convex set S ⊂ Rn and let
x∗ be an interior point of S. Assume that f (·) is C 1 in a ball around x∗ .

1. If f (·) is concave in S, then x∗ is a (global) maximum point for f (·) in S if and

only if x∗ is a stationary point for f (·).

2. f (·) is convex in S, then x∗ is a (global) minimum point for f (·) in S if and only
if x∗ is a stationary point for f (·).

Proof of Theorem 6.2: We focus on the first part of the theorem. The second part
follows once we take into account that −f is concave. (=⇒) This follows from Theorem
6.1 above. (⇐=) Suppose that x∗ is a stationary point for f (·) and that f (·) is concave.
Recall the inequality in Theorem 5.6 (First-order characterization of concave functions).
For any x ∈ S,

f (x) − f (x∗ ) ≤ ∇f (x∗ ) · (x − x∗ ) = 0 (∵ ∇f (x∗ ) = 0)

Thus, we have f (x) ≤ f (x∗ ) for any x ∈ S as desired. 

6.1.2 Envelope Theorems for Unconstrained Maxima

Consider an objective function with a parameter vector r of the form f (x, r) =
f (x1 , . . . , xn , r1 , . . . , rk ), where x ∈ S ⊂ Rn and r ∈ Rk . For each fixed r, suppose we
have found the maximum of f (x, r) when x varies in S. The maximum value of f (x, r)
usually depends on r. We denote this value by f ∗ (r) and call f ∗ the value function.
Thus,

f ∗ (r) = max f (x, r)

x∈S

The vector x that maximizes f (x, r) depends on r and is therefore denoted by x∗ (r).
Then, f ∗ (r) = f (x∗ (r), r).

Theorem 6.3 (Envelope Theorem) In the maximization problem maxx∈S f (x, r), where
S ⊂ Rn and r ∈ Rk , suppose that there is a maximum point x∗ (r) ∈ S for every
r ∈ Bδ (r∗ ) with some δ > 0. Furthermore, assume that the mappings r → f (x∗ (r∗ ), r)
and r → f ∗ (r) are differentiable at r ∗ . Then
⎡ ⎤
∂f (x, r)/∂r1
⎢ ∂f (x, r)/∂r2 ⎥
⎢ ⎥
∇r f ∗ (r∗ ) = ⎢ .. ⎥
⎣ . ⎦
∂f (x, r)/∂rk x=x∗ (r∗ ),r=r ∗

78
 There are two effects of r on the value function f ∗ through both directly and indirectly
x∗ (r).
The Envelope theorem says that we can ignore the indirect effects.

Proof of Envelope Theorem: Define the function

ϕ(r) = f (x∗ (r∗ ), r) − f ∗ (r).

Because x∗ (r∗ ) is a maximum point of f (x, r) when r = r ∗ , one has ϕ(r∗ ) = 0 and
ϕ(r) ≤ 0 for all r ∈ Bδ (r∗ ). Since ϕ(r∗ ) is a maximum, the following first order condition
is satisfied (because of Theorem 6.1).

 ∂ϕ(r) 

∇r ϕ(r) r=r∗ = 0 ⇐⇒ = 0 ∀j = 1, . . . , k
∂rj r=r∗

That is,
 
∂ϕ(r) ∂f (x∗ (r∗ ), r)  ∂f ∗ (r) 
∂rj
=
∂rj  ∗ − ∂rj  ∗ = 0 ∀j = 1, . . . , k 
r=r r=r

6.1.3 Local Extreme Points

The point x∗ is a local maximum point of f (·) in S if there exists an ε > 0 such that
f (x) ≤ f (x∗ ) for all x ∈ Bε (x∗ ) ∩ S. If x∗ is the unique local maximum point for f (·),
then it is a strict local maximum point for f (·) in S. A (strict) local minimum point is
defined in the obvious way, and it should be clear what we mean by local maximum and
minimum values, local extreme points, and local extreme values. A stationary point x∗ of
f (·) that is neither a local maximum point nor a local minimum point is called a saddle
point of f (·).

Before stating the next result, recall the n leading principal minors of the Hessian
matrix D 2 f (x):
 
 f11 (x) f12 (x) ··· f1k (x) 
 
 f21 (x) f22 (x) ··· f2n (x) 
2  
|D(k) f (x)| =  .. .. .. ..  , k = 1, . . . , n
 . . . . 
 
 fk1 (x) fk2 (x) ··· fkk (x) 

Theorem 6.4 (Sufficient Conditions for Local Extreme Points) Suppose that f (x) =
f (x1 , . . . , xn ) is defined on a set S ⊂ Rn and that x∗ is an interior stationary point. As-
sume also that f (·) is C 2 in an open ball around x∗ . Then,

1. D 2 f (x∗ ) is positive definite =⇒ x∗ is a local minimum point.

2. D 2 f (x∗ ) is negative definite =⇒ x∗ is a local maximum point.

79
 Proof of Theorem 6.4: We only focus on the first part of the theorem. We should
be able to prove the second part of the proof by replacing f (·) with −f (·). Since each
fij (x) is continuous in x (because f (·) is C 2 ), the determinant is a continuous function
2 f (x∗ )| > 0 for all k, it is possible to find a ball B (x∗ ) with ε > 0
of x. Therefore, if |D(k) ε
2 f (x)| > 0 for all x ∈ B (x∗ ) and all k = 1, . . . , n. By Theorem 4.6,
so small that |D(k) ε
the corresponding quadratic form is positive definite for all x ∈ Bε (x∗ ). It follows from
Theorem 5.5 that f (·) is strictly convex in Bε (x∗ ). Then, Theorem 6.2 shows that the
stationary point x∗ is a maximum point for f in Bε (x∗ ). Hence, x∗ is a local minimum
point for f (·). 

Lemma 6.1 If x∗ is an interior stationary point of f (·) such that |D2 f (x∗ )|
= 0 and
D 2 f (x∗ ) is neither positive definite nor negative definite, then x∗ is a saddle point.

6.1.4 Necessary Conditions for Local Extreme Points

To study the behavior of f (·) in an arbitrary fixed vector in Rn with length 1, so h = 1.
The function g(·) describes the behavior of f (·) along the straight line through x∗ parallel
to the vector h ∈ Rn .

g(t) = f (x∗ + th) = f (x∗1 + th1 , . . . , x∗n + thn )

We have the following characterization of local extreme points.

Theorem 6.5 (Necessary Conditions for Local Extreme Points) Suppose that f (x) =
f (x1 , . . . , xn ) is defined on a set S ⊂ Rn , and x∗ is an interior stationary point in S.
Assume that f is C 2 in a ball around x∗ . Then,

1. x∗ is a local minimum point =⇒ D2 f (x∗ ) is positive semidefinite.

2. x∗ is a local maximum point =⇒ D2 f (x∗ ) is negative semidefinite.

Proof of Theorem 6.5: Suppose that x∗ is an interior local maximum point for
f (·). Then, if ε > 0 is small enough, Bε (x∗ ) ⊂ S, and f (x) ≤ f (x∗ ) for all x ∈ Bε (x∗ ).
If t ∈ (−ε, ε), then x∗ + th ∈ Bε (x∗ ) because (x∗ + th) − x∗  = th = |t| < ε. Then,
for all t ∈ (−ε, ε), we have

f (x∗ + th) ≤ f (x∗ ) ⇐⇒ g(t) ≤ g(0).

Thus, the function g(·) has an interior maximum at t = 0. Using the chain rule, we
obtain
n


g (t) = fi (x∗ + th)hi
i=1
n
n

g (t) = fij (x∗ + th)hi hj
i=1 j=1

80
 
The condition g (0) ≤ 0 yields
n
n
fij (x∗ + th)hi hj ≤ 0 ∀h = (h1 , . . . , hn ) with h = 1
i=1 j=1

This implies that the Hessian matrix D 2 f (x∗ ) is negative semidefinite. Theorem 4.6
shows that this is equivalent to checking all principal minors. The same argument can
be used to establish the necessary condition for x∗ to be a local minimum point for f (·).


Exercise 6.1 Find the local extreme values and classify the stationary points as maxima,
minima, or neither.

1. f (x1 , x2 ) = 2x1 − x21 − x22 .

2. f (x1 , x2 ) = x21 + 2x22 − 4x2 .

3. f (x1 , x2 ) = x31 − x22 + 2x2 .

4. f (x1 , x2 ) = 4x1 + 2x2 − x21 + x1 x2 − x22 .

5. f (x1 , x2 ) = x31 − 6x1 x2 + x32

6.2 Constrained Optimization

6.2.1 Equality Constraints: The Lagrange Problem
A general maximization problem with equality constraints is of the form

max f (x1 , . . . , xn ) subject to gj (x) = 0 ∀j = 1, . . . , m (m < n) (∗)

x=(x1 ,... ,xn )

Define the Lagrangian,

L(x) = f (x) − λ1 g1 (x) − · · · − λm gm (x)

where λ1 , . . . , λm are called Lagrange multipliers. The necessary first-order conditions

for optimality are then:
m
∂L(x) ∂f (x) ∂gj (x)
∇L(x) = ∇f (x) − λDg(x) = 0 ⇐⇒ = − λj = 0, ∀i = 1, . . . , n (∗∗)
∂xi ∂xi ∂xi
j=1

Theorem 6.6 (N&S Conditions for Extreme Points with Equality Constraints)
The following establishes the necessary and sufficient conditions for the Lagrangian method

81
 1. (Necessity) Suppose that the functions f and g1 , . . . , gm are defined on a set S
in Rn and x∗ = (x∗1 , . . . , x∗n ) is an interior point of S that solves the maximization
problem (∗). Assume further that f and g1 , . . . , gm are C 1 in a ball around x∗ , and
that the Jacobian matrix,
⎛ ⎞
∂g 1 (x∗ ) ∂g 1 (x∗ )
· · ·
⎜ ∂x. 1 ∂xn ⎟
Dg(x∗ ) = ⎜ ⎝ .
.
..
.
..
.
⎟
⎠
   m ∗ m ∗
∂g (x ) ∂g (x )
m×n
∂x1 · · · ∂xn

has rank m. Then, there exist unique numbers λ1 , . . . , λm such that the first-order
conditions (∗∗) are valid.

2. (Sufficiency) If there exist numbers λ1 , . . . , λm and a feasible x∗ which together

satisfy the first-order conditions (∗∗), and if the Lagrangian L(x) is concave in x,
then x∗ solves the maximization problem (∗).

Proof of Theorem 6.6: (Necessity) The proof for the necessity part consists of
three steps.

Step 1: Construction of the unconstrained maximization problem

Since m × n matrix Dg(x∗ ) is assumed to have rank m, there exists a invertible (non-
singular) m × m submatrix. After renumbering the variables, if necessary, we can assume
that it consists of the first m columns. By the implicit function theorem (Theorem 5.17),
the m constraints, g1 , . . . , gm define x1 , . . . , xm as C 1 functions of x̃ = (xm+1 , . . . , xn )
in some open ball B around x∗ , i.e., Bε (x∗ ) with ε > 0 sufficiently small, so we can write

xj = hj (xm+1 , . . . , xn ) = hj (x̃), j = 1, . . . , m

Then, f (x1 , . . . , xn ) reduces to a composite function

ξ(x̃) = f (h1 (x̃), . . . , hm (x̃), x̃)

of x̃ only. Now, the maximization problem with equality constraints is translated into
the unconstrained maximization problem

max ξ(x̃)
x̃∈Bε (x∗ )

Since x∗ is a local extreme point for f subject to the given constraints, ξ must have an
unconstrained local extreme point at x̃∗ = (x∗m+1 , . . . , x∗n ). Hence, the partial derivatives
of ξ with respect to xm+1 , . . . , xn must be 0:


∂ξ(x∗ ) ∂f (x∗ ) ∂h1 ∂f (x∗ ) ∂hm ∂f (x∗ )
= + ··· + + = 0 (1)
∂xk ∂x1 ∂xk ∂xm ∂xk ∂xk
     
indirect effects direct effects

82
where k = m + 1, . . . , n.

Step 2: To express ∂h1 /∂xk , . . . , ∂hm /∂xk in terms of ∇g j (x)

gj (h1 (x̃), . . . , hm (x̃), x̃) = 0, j = 1, . . . , m

for all x̃ ∈ B. Differentiating this with respect to xk gives

m
∂gj ∂hs ∂gj
+ = 0,
∂xs ∂xk ∂xk
s=1

for k = m+1, . . . , n and j = 1, . . . , m. In particular, this is valid at x∗ . Now, multiplying

each of the m equations above by a scalar λj , then adding these equations over j, we
obtain
m
m  m
∂gj ∂hs ∂gj
λj + λj = 0, (2)
s=1
∂xs ∂xk ∂xk
j=1 j=1

where k = m + 1, . . . , n. Next, subtract (2) from (1),

⎛ ⎞
m m m
⎝ ∂f (x ) −
∗ j
∂g ⎠ ∂hs ∂f (x∗ ) ∂gj
λj + − λj =0
∂xs ∂xs ∂xk ∂xk ∂xk
s=1 j=1 j=1

where k = m + 1, . . . , n. These equations are valid for all choices of λ1 , . . . , λm when

the partial derivatives are evaluated at x∗ . Suppose that we can prove the existence of
numbers λ1 , . . . , λm such that
m
∂f (x∗ ) ∂gj
− λj = 0, s = 1, . . . , m (3)
∂xs ∂xs
j=1

If (3) is satisfied, we also have

m
∂f (x∗ ) ∂gj
− λj ∀k = m + 1, . . . , n
∂xk ∂xk
j=1

Thus, the first-order necessary conditions for the Lagrangian are also satisfied.

Step 3: Existence of λ1 , . . . , λm satisfying (3)

83
 We rewrite the system (3) as
∂g1 ∂g2 ∂gm ∂f
λ1 + λ2 · · · λm =
∂x1 ∂x1 ∂x1 ∂x1
∂g1 ∂g2 ∂gm ∂f
λ1 + λ2 · · · λm = ⇐⇒ Dg(x∗ ) 
λ = ∇f (x∗ )
∂x2 ∂x2 ∂x2 ∂x1      
m×n n×1 m×1
··························· ··· ···
∂g1 ∂g2 ∂gm ∂f
λ1 + λ2 · · · λm =
∂xm ∂xm ∂xm ∂xm
Here the coefficient matrix above is invertible (nonsingular), it has a unique solution
λ1 , . . . , λm . This completes the proof.

(Sufficiency) Suppose that the Lagrangian L(x) is concave. The first-order neces-
sary conditions imply that the Lagrangian is stationary at x∗ . Then by Theorem 6.2
(sufficiency for global maximum),
m
m

∗ ∗ j ∗
L(x ) = f (x ) − λj g (x ) ≥ f (x) − λj gj (x) = L(x) ∀x ∈ S
j=1 j=1

But for all feasible x, we have gj (x) = 0 and of course, gj (x∗ ) = 0 for all j = 1, . . . , m.
Hence, this implies that f (x∗ ) ≥ f (x). Thus, x∗ solves the maximization problem (∗). 

6.2.2 Lagrange Multipliers as Shadow Prices

The optimal values of x1 , . . . , xn in the maximization problem (∗) will depend upon the
parameter vector r = (r1 , . . . , rk ) ∈ Rk , in general. If x∗ (r) = (x∗1 (r), . . . , x∗n (r)) denotes
the vector of optimal values of the choice variables, then the corresponding value
f ∗ (r) = f (x∗1 (r), . . . , x∗n (r))
of f (·) is called the (optimal) value function for the maximization problem (∗). The values
of the Lagrange multipliers  will also depend on r; we write λj = λj (r) for j = 1, . . . , m.
Let L(x, r) = f (x, r) − m j=1 j λ g j (x, r) be the Lagrangian. Under certain conditions, we

have


∂f ∗ (r) ∂L(x, r)
= , i = 1, . . . , k
∂ri ∂ri x=x∗ (r)

6.2.3 Tangent Hyperplane

A set of equality constraints in Rn
g1 (x) = 0
g2 (x) = 0
.. ..
. .
gm (x) = 0

84
defines a subset of Rn which is best viewed as a hypersurface. If, as I assume in this
section, the functions gj , j = 1, . . . , m belong to C 1 , the surface defined by them is said
to be smooth. We introduce the tangent hyperplane M below:
M = {y ∈ Rn |Dg(x∗ )y = 0}
Note that the tangent hyperplane is a subspace of Rn .

Definition 6.1 A point x∗ satisfying the constraint g(x∗ ) = 0 is said to be a regu-

lar point of the constraint if the gradient vectors ∇g 1 (x∗ ), . . . , ∇g m (x∗ ) are linearly
independent. That is, Rank(Dg(x∗ )) = m.

6.2.4 Local First-Order Necessary Conditions

Lemma 6.2 Let x∗ be a regular point of the constraint g(x) = 0 and a local extreme
point of f (·) subject to these constraints. Then, for any y ∈ Rn ,
∇f (x∗ )y = 0 whenever Dg(x∗ )y = 0.
Proof of Lemma 6.2: Let y = (y1 , . . . , yn ) with y = 1. Let x(t) = x∗ + ty be any
smooth curve on the constraint surface g(x(t)) = 0 passing through x∗ with derivative

x (t)|t=0 = y at x(0) = x∗ . There exists some ε > 0 such that g(x(t)) = 0 for any
t ∈ (−ε, ε).

Since x∗ is a regular point, the tangent hyperplane is identical with the set of y’s
satisfying ∇g(x∗ )y = 0. Then, since x∗ is a constrained local extreme point of f (·), we
have

d 
f (x(t))

= 0 =⇒ ∇f (x∗ )x (0) = 0,
dt t=0
equivalently, ∇f (x∗ )y = 0. 

The above lemma says that ∇f (x∗ ) is orthogonal to the tangent hyperplane.

6.2.5 Second-Order Necessary and Sufficient Conditions for Local Ex-

treme Points
Theorem 6.7 (Necessity for Local Maximum) Suppose that x∗ is a local maximum
of f (·) subject to g(x) = 0 and that x∗ is a regular point of these constraints. Then there
is a λ ∈ Rn such that
∇f (x∗ ) − λDg(x∗ ) = 0.
If we denote by M the tangent hyperplane M = {h ∈ Rn |Dg(x∗ )h = 0}, then, the matrix
D 2 L(x∗ ) = D2 f (x∗ ) − λD2 g(x∗ )
is negative semidefinite on M , that is,
hT D2 L(x∗ )h ≤ 0 ∀h ∈ M

85
 Proof of Theorem 6.7: The first part follows from Theorem 6.6. We only focus
on the second part. Let h = (h1 , . . . , hn ) ∈ M with h = 1. Let x(t) = x∗ + th be
any smooth curve on the constant surface g(x(t)) = 0 passing through x∗ with derivative

x (0) = h at x(0) = x∗ . Suppose that x∗ is an interior local maximum point for f subject
to g(x) = 0. Then, if ε > 0 is small enough,

L(x∗ + th) ≤ L(x∗ ) ⇐⇒ f (x∗ + th) − λg(x∗ + th) ≤ f (x∗ ) − λg(x∗ )

for all t ∈ (−ε, ε) because (x∗ + th) − x∗  = th = |t| < ε. Define the function
ϕ(t) = L((x∗ + th). Then, for all t ∈ (−ε, ε), we have

L(x∗ + th) ≤ L(x∗ ) ⇐⇒ ϕ(t) ≤ ϕ(0).

Thus, the function ϕ has an interior maximum at t = 0. Using the chain rule (Theorem
5.15), we obtain

ϕ (t) = ∇L(x∗ + th)h = ∇f (x∗ + th)h − λDg(x∗ + th)h

ϕ (0) = ∇L(x∗ )h = ∇f (x∗ )h − λDg(x∗ )h = 0

because h ∈ M so that ∇f (x∗ )h = 0 and Dg(x∗ )h = 0. Furthermore,


ϕ (t) = hT D2 L(x∗ + th)h

The hypothesis that ϕ has an interior local maximum at t = 0 means ϕ (0) ≤ 0. Thus,
n

hT D2 L(x∗ )h ≤ 0 ⇐⇒ Lij (x∗ )hi hj ≤ 0
i=1

This implies that the Hessian matrix D 2 L(x∗ ) is negative semidefinite on M . 

Theorem 6.8 (Sufficiency for Local Maximum) Suppose there is a point x∗ ∈ Rn

satisfying g(x∗ ) = 0, and a λ ∈ Rm such that
⎛ ⎞
0
∗  ∗ ⎜ .. ⎟
∇f (x ) − λ Dg(x ) = 
0 = ⎝ . ⎠.
n×1 0

Suppose also that the matrix D 2 L(x∗ ) = D2 f (x∗ ) + λD2 g(x∗ ) is negative definite on
M = {y ∈ Rn |Dg(x∗ )y = 0}, that is, for y ∈ M with y
= 0, y T D2 L(x∗ )y < 0. Then, x∗
is a strict local maximum of f (·) subject to g(x) = 0.
Proof of Theorem 6.8: The first part follows from Theorem 6.7. Define the La-
grangian as follows.

L(x) = f (x) − λg(x)

Differentiating this with respect to x, and evaluating it at x∗ , we obtain

∇L(x∗ ) = ∇f (x∗ ) − λDg(x∗ ) = 0

86
This implies that ∇L(x∗ )y = 0 for any y ∈ Rn . By our hypothesis, D2 L(x∗ ) is negative
definite on M , and therefore, x∗ is a local maximum point of L(x) from Theorem 6.4.
This implies that x∗ is a local maximum of f (·) subject to g(x) = 0. 

Exercise 6.2 Solve the problem

max{x + 4y + z} subject to x2 + y 2 + z 2 = 216 and x + 2y + 3z = 0

Exercise 6.3 Consider the problem (assuming m ≥ 4).

1 1
max U (x1 , x2 ) = ln(1 + x1 ) + ln(1 + x2 ) subject to 2x1 + 3x2 = m
2 4
Answer the following questions.

1. Let x∗1 (m) and x∗2 (m) denote the values of x1 and x2 that solve the above maximiza-
tion problem. Find these functions and the corresponding Lagrangian multiplier.

2. The optimal value U ∗ of U (x1 , x2 ) is a function of m. Find an explicit expression

for U ∗ (m), and show that dU ∗ /dm = λ.

6.2.6 Envelope Result for Lagrange Problems

In economic optimization problems, the objective function as well as the constraint func-
tions (such as the budget set) will often depend on parameters. These parameters are
held constant when optimizing (remember the price taking behavior assumption), but
can vary with the economic situation. We might want to know what happens to the
optimal value function when the parameter change.

Consider the following general Lagrange problem.

max f (x, r) subject to gj (x, r) = 0, j = 1, . . . , m.

x∈S

where r = (r1 , . . . , rk ) is a vector of parameters. The values of x1 , . . . , xn that solve

the maximization problem will be functions of r. If we denote them by x∗1 (r), . . . , x∗n (r),
then

f ∗ (r) = f (x∗1 (r), . . . , x∗n (r))

is called the value function. Suppose that λi = λi (r) for all i = 1, . . . , m are the Lagrange
multipliers in the first-order conditions for the maximization problem and let

L(x, r) = f (x, r) + λ · g(x, r)

be the Lagrangian. Here λ = (λ1 , . . . , λm ) and g(x, r) = (g1 (x, r), . . . , gm (x, r)).

87
6.3 Inequality Constraints: Nonlinear Programming
⎧ 1
⎪
⎪ g (x1 , . . . , xn ) ≤ 0
⎪
⎨ g2 (x1 , . . . , xn ) ≤ 0
max f (x) subject to ..
x∈S ⎪
⎪ .
⎪
⎩ m
g (x1 , . . . , xn ) ≤ 0

A vector x = (x1 , . . . , xn ) that satisfies all the constraints is called feasible. The
set of all feasible vectors is said to be the feasible set. We assume that f (·) and all the
gj functions are C 1 . In the case of equality constraint, the number of constraints were
assumed to be strictly less than the number of variables. This is not necessary for the
case of inequality constraints. An inequality constraint gj (x) ≤ 0 is said to be active
(binding) at x if gj (x) = 0 and inactive (non-binding) at x if gj (x) < 0.

Note that minimizing f (x) is equivalent to maximizing −f (x). Moreover, an inequal-

ity constraint of the form gj (x) ≥ 0 can be rewritten as −g j (x) ≤ 0. In this way, most
constrained optimization problem can be expressed as the above form.

We define the Lagrangian exactly as before.

m

L(x) = f (x) − λ · g(x) = f (x) − λj gj (x),
j=1

where λ = (λ1 , . . . , λm ) ∈ Rm are the Lagrangian multipliers. Again the first-order

partial derivatives of the Lagrangian are equated to 0:
m
∂L(x) ∂f (x) ∂gj (x)
∇L(x) = ∇f (x) − λDg(x) = 0 ⇐⇒ = − λj = 0, ∀i = 1, . . . , n (∗)
∂xi ∂xi ∂xi
j=1

In addition, we introduce the complementary slackness conditions. For all j = 1, . . . , m,

λj ≥ 0 and λj = 0 if gj (x) < 0 (∗∗)

An alternative formulation of this condition is that for any j = 1, . . . , m,

λj ≥ 0 and λj gj (x) = 0

In particular, if λj > 0, we must have gj (x) = 0. However, it is perfectly possible to have

both λj = 0 and gj (x) = 0.

Conditions (∗) and (∗∗) are often called the Kuhn-Tucker conditions. They are
(essentially but not quite) necessary conditions for a feasible vector to solve the maxi-
mization problem. In general, they are definitely not sufficient on their own. Suppose
one can find a point x∗ at which f (·) is stationary and gj (x∗ ) < 0 for all j = 1, . . . , m.
Then, the Kuhn-Tucker conditions will automatically be satisfied by x∗ together with all
the Lagrangian multipliers λj = 0 for all j = 1, . . . , m.

88
Theorem 6.9 (Sufficiency for the Kuhn-Tucker Conditions I) Consider the max-
imization problem and suppose that x∗ is feasible and satisfies conditions (∗) and (∗∗).
If the Lagrangian L(x) = f (x) − λ · g(x) (with the λ values obtained from the recipe) is
concave, then x∗ is optimal.
Proof of Theorem 6.9: This is very much the same as the sufficiency part of
the Lagrangian problem in Theorem 6.6. Since L(x) is concave by assumption and
∇L(x∗ ) = 0 from (∗), by Theorem 6.2, x∗ is a global maximum point of L(x). Hence,
for all x ∈ S,
m
m

f (x∗ ) − λj gj (x∗ ) ≥ f (x) − λj gj (x)
j=1 j=1

Rearranging gives the equivalent inequality

m
!
f (x∗ ) − f (x) ≥ λj gj (x∗ ) − gj (x) .
j=1

Thus, it suffices to show

m
!
λj gj (x∗ ) − gj (x) ≥ 0
j=1

for all feasible x, because this will imply that x∗ solves the maximization problem. Sup-
pose that gj (x∗ ) < 0. Then (∗∗) shows that λj = 0. Suppose that gj (x∗ ) = 0, we have
λj (gj (x∗ ) − gj (x)) j
m = −λjjg (x) ≥ 0 because
! x is feasible, i.e., gj (x) ≤ 0 and λj ≥ 0.
∗ j
Then, we have j=1 λj g (x ) − g (x) ≥ 0 as desired. 

Theorem 6.10 (Sufficiency for the Kuhn-Tucker Conditions II) Consider the max-
imization problem and suppose that x∗ is feasible and satisfies conditions (∗) and (∗∗).
If f (·) is concave and each λj gj (x) (with the λ values obtained from the recipe) is qua-
siconvex, then x∗ is optimal.
Proof of Theorem 6.10: We want to show that f (x) − f (x∗ ) ≤ 0 for all feasible
x. Since f (·) is concave, then according to Theorem 5.6 (First-order characterization of
concavity of f (·)),
m

f (x) − f (x∗ ) ≤ ∇f (x∗ ) · (x − x∗ ) 
= λj ∇g j (x∗ ) · (x − x∗ )
(∗) j=1

where we use the first order condition (∗). It therefore suffices to show that for all
j = 1, . . . , m, and all feasible x,

λj ∇g j (x∗ ) · (x − x∗ ) ≤ 0.

The above inequality is satisfied for those j such that gj (x∗ ) < 0, because then λj = 0
from the complementary slackness condition (∗∗). For those j such that gj (x∗ ) = 0, we

89
have gj (x) ≤ gj (x∗ ) (because x is feasible), and hence −λj gj (x) ≥ −λj gj (x∗ ) because
λj ≥ 0. Since the function −λj gj (x) is quasiconcave (because λj gj (x) is quasiconvex),
it follows from Theorem 5.10 (a characterization of quasiconcavity) that ∇(−λj gj (x∗ )) ·
(x − x∗ ) ≥ 0, and thus, λj ∇g j (x∗ ) · (x − x∗ ) ≤ 0. 

Exercise 6.4 Reformulate the problem

min 4 ln(x2 + 2) + y 2 subject to x2 + y ≥ 2, x ≥ 1

as a standard Kuhn-Tucker maximization problem and write down the necessary Kuhn-
Tucker conditions. Moreover, find the solution of the problem (Take it for granted that
there is a solution).

6.4 Properties of the Value Function

Consider the standard nonlinear programming problem below.

max f (x) subject to gj (x) ≤ bj ∀j = 1, . . . , m

The optimal value of the objective f (x) obviously depends upon b ∈ Rm which is a
parameter vector in the constraint set. The function defined by
 
f ∗ (b) = max f (x)gj (x) ≤ bj , j = 1, . . . , m

assigns to each b = (b1 , . . . , bk ) the optimal value f ∗ (b) of f (·). It is called the value
function for the problem. Let the optimal choice for x in the constrained optimization
problem be denoted by x∗ (b), and assume that it is unique. Let λj (b) for j = 1, . . . , m
be the corresponding Lagrange multipliers. Then, if ∂f ∗ (b)/∂bj exists,
∂f ∗ (b)
= λj (b) ∀j = 1, . . . , m
∂bj

The value function f ∗ is not necessarily C 1 . The next proposition characterizes a geo-
metric structure of the value function.

Proposition 6.1 If f (x) is concave and gj (x) is convex for each j = 1, . . . , m, then
f ∗ (b) is concave.
 
Proof of Proposition 6.1: Suppose that b and b are two arbitrary parame-
   
ter vectors in the constraint set, and let f ∗ (b ) = f (x∗ (b ), f ∗ (b ) = f (x∗ (b ), with
   
gj (x∗ (b )) ≤ bj , gj (x∗ (b )) ≤ bj for j = 1, . . . , m. Let α ∈ [0, 1]. Corresponding to the
   
vector αb + (1 − α)b , there exists an optimal solution x∗ (αb + (1 − α)b ), and
   
f ∗ (αb + (1 − α)b ) = f (x∗ (αb + (1 − t)b ))
 
Define xα = αx∗ (b ) + (1 − α)x∗ (b ). Then, convexity of gj for j = 1, . . . , m implies that
   
gj (xα ) ≤ αgj (x∗ (b )) + (1 − α)gj (x∗ (b )) ≤ αbj + (1 − α)bj

90
  
Thus, xα is feasible in the problem with parameter αb + (1 − α)b , and in that problem,
 
x∗ (αb + (1 − α)b ) is optimal. It follows that
   
f (xα ) ≤ f (x∗ (αb + (1 − α)b )) = f ∗ (αb + (1 − α)b )

On the other hand, concavity of f implies that

   
f (xα ) ≥ αf (x∗ (b )) + (1 − α)f (x∗ (b )) = αf ∗ (b ) + (1 − α)f (b )

In sum,
   
f ∗ (αb + (1 − α)b ) ≥ αf ∗ (b ) + (1 − α)f (b )

This shows that f ∗ (b) is concave. 

6.5 Constraint Qualifications

Consider the maximization problem below.

max f (x) subject to gj (x) ≤ 0, j = 1, . . . , m

Definition 6.2 The constrained maximization problem satisfies the constraint qual-
ification if the gradient vectors ∇g j (x∗ ) (1 ≤ j ≤ m) corresponding those constraints
that are active (binding) at x∗ , are linearly independent.
An alternative formulation of this condition is: Delete all rows in the Jacobian matrix
Dg(x∗ ) that correspond to constraints that are inactive (not binding) at x∗ . Then, the
remaining matrix should have rank equal to the number of rows.

Theorem 6.11 (Kuhn-Tucker Necessary Conditions) Suppose that x∗ = (x∗1 , . . . , x∗n )

solves the constrained maximization problem where f (·) and g1 (·), . . . , gm (·) are C 1 func-
tions. Suppose furthermore that the maximization problem satisfies the constraint qual-
ification. Then, there exist unique numbers λ1 , . . . , λm such that the Kuhn-Tucker con-
ditions (∗) and (∗∗) hold at x = x∗ .
Proof of Kuhn-Tucker Necessary Conditions: We assume the following.

1. x∗ ∈ Rn maximizes f on the constraint set gj (x) ≤ 0 for all j = 1, . . . , m

2. only g1 , . . . , gk are binding at x∗ , where k ≤ m.

3. the k × n Jacobian matrix Dgk (x∗ ) has maximal rank k. That is,
⎛ ⎞
∂g1 (x∗ )/∂x1 · · · ∂g1 (x∗ )/∂xn
⎜ .. .. ⎟
k = Rank [Dgk (x∗ )] = Rank ⎝ .
..
. . ⎠.
k ∗ k ∗
∂g (x )/∂x1 · · · ∂g (x )/∂xn

91
 The proof consists of two steps.

Step 1: ∇L(x, λ) = 0 and λg(x) = 0

Since each gj (·) is a continuous function, there is a open ball Bε (x∗ ) such that gj (x) <
0 for all x ∈ Bε (x∗ ) and for j = k + 1, . . . , m. We will work in the open ball Bε (x∗ ) for
the rest of proof.

Note that x∗ maximizes f (·) in Bε (x∗ ) over the constraint set that gj (x) = 0 for
j = 1, . . . , k. By assumption, Theorem 6.7 (Necessity for Optimization with Equality
Constraints) applies and therefore, there exist μ∗1 , . . . , μ∗k such that

∇L̂(x∗ , μ∗ ) = 0 and gj (x∗ ) = 0 ∀j = 1, . . . , k


where L̂(x, μ) ≡ f (x) − kj=1 μj gj (x) as the restricted Lagrangian.

Consider the usual Lagrangian

m

L(x, λ1 , . . . , λm ) ≡ f (x) − λk gj (x).
j=1

Let λ∗i = μ∗i for i = 1, . . . , k and λ∗i = 0 for j = k + 1, . . . , m. Then, we see that (x∗ , λ∗ )
is a solution of the n + m equations in n + m unknowns:
∂L ∗ ∗
(x , λ ) = 0 ∀i = 1, . . . , n
∂xi
λ∗j gj (x∗ ) = 0 ∀j = 1, . . . , m

Step 2: λj ≥ 0 for all j

There is a C 1 curve x(t) defined for t ∈ [0, ε) such that x(0) = x∗ and, for all t ∈ [0, ε),

g1 (x(t)) = −t and gj (x(t)) = 0 for j = 2, . . . , k

By the implicit function theorem (Theorem 5.17), we can still solve the constrained
optimization problem in Bε (x∗ ) even if we slightly perturb the constraint set. Let h =

x (0). Using the chain rule (Theorem 5.15), we conclude that

nablag 1 (x∗ )h = −1, ∇g j (x∗ )h = 0 ∀j = 2, . . . , k

Since x(t) lies in the constraint set for all t and x∗ maximizes f (·) in the constraint set,
f (·) must be nonincreasing along x(t). Therefore,

d 
f (x(t)) = ∇f (x∗ )h ≤ 0
dt t=0

92
By our first-order conditions, we execute a series of computations.

0 = ∇L(x∗ )y
k

= ∇f (x∗ )h − λj ∇g j (x∗ )y
j=1

= ∇f (x )h − λ1 ∇g 1 (x∗ )h
∗

= ∇f (x∗ )h + λ1

Since ∇f (x∗ )h ≤ 0, we conclude that λ1 ≥ 0. A similar argument shows that λj ≥ 0 for

j = 1, . . . , k. This completes the proof. 

Theorem 6.12 (Kuhn-Tucker N & S Conditions) Assume that a feasible vector x∗

and a set of multipliers λ1 , . . . , λm satisfy the Kuhn-Tucker necessary conditions (∗) and
(∗∗) for the constrained maximization problem. Define J = {j|gj (x∗ ) = 0}, the set
of active (binding) constraints, and assume that λj > 0 for all j ∈ J Consider the
Lagrangian problem

max f (x) subject to gj (x) = 0 ∀j ∈ J

Then, x∗ satisfies

∇L̃(x∗ ) = ∇f (x∗ ) − λj ∇g j (x∗ ) = 0
j∈J

for the given multipliers λj for j ∈ J. If D2 L̃(x∗ ) is negative definite on M , then x∗ is

a strict local maximum point for the original constrained maximization problem. Here
 
M = h ∈ Rn  ∇g j (x∗ )h = 0 ∀j ∈ J

Proof of Kuhn-Tucker N & S Conditions: Suppose, on the contrary, that x∗

is not a local maximum point of the constrained optimization problem. Then, we can
consider {yk } as a sequence of feasible points converging to x∗ such that f (yk ) ≥ f (x∗ )
for each k. More specifically, for each k, define yk = x∗ + εk hk with hk  = 1 and εk > 0.
We may assume that εk → 0 and hk → h∗ as k → ∞. Using the linear approximation
through differentiability,

f (yk ) ≈ f (x∗ ) + ∇f (x∗ ) · (yk − y ∗ ) = f (x∗ ) + εk ∇f (x∗ ) · hk

for k large enough. Letting k → ∞, because of linearity of ∇f (x∗ )hk in hk (continuity

follows), we must have ∇f (x∗ )h∗ ≥ 0 from f (yk ) ≥ f (x∗ ). Also for each binding (active)
constraint gj , we have

gj (yk ) ≤ gj (x∗ )

Again, using the linear approximation through differentiability,

gj (yk ) ≈ gj (x∗ ) + ∇g j (x∗ ) · (yk − x∗ ) = gj (x∗ ) + εk ∇g j (x∗ ) · hk

93
for k large enough. Then, we must have Dgj (x∗ )h∗ ≤ 0 because Dgj (x∗ )hk is a linear
continuous in hk and gj (yk ) ≤ gj (y ∗ ) for each k.

If ∇g j (x∗ )h∗ = 0 for all j ∈ J, then the proof goes through just as in Theorem 6.8.
Therefore, there must exists at least one j ∈ J such that ∇g j (x∗ )h∗ < 0. Then, we
obtain

∇f (x∗ )h∗ − λj Dgj (x∗ )h∗ > 0 because λj > 0 for all j ∈ J
j∈J
⎡ ⎤

⎣∇f (x∗ ) − λj Dgj (x∗ )⎦ h∗ > 0
j∈J
  
=0

This, however, contradicts our fulfilled condition that ∇f (x∗ ) − j∈J λj ∇g j (x∗ ) = 0.
We complete the proof. 

Exercise 6.5 Consider the following constrained maximization problem.

max f (x, y) = x subject to g(x, y) = x3 + y 2 = 0.
Show that this problem does not satisfy the constraint qualification.

6.6 Nonnegativity Constraints

Consider the nonlinear programming problem with nonnegativity constraints:
max f (x) subject to gj (x) ≤ 0 ∀j = 1, . . . , m and xi ≥ 0 for all i = 1, . . . , n
We introduce n new constraints in addition to the m original ones:
gm+1 (x) = −x1 ≤ 0
gm+2 (x) = −x2 ≤ 0
.. .. ..
. . .
gm+n (x) = −xn ≤ 0
We introduce the Lagrangian multipliers μ1 , . . . , μn to go with the new constraints and
form the extended Lagrangian.
m
n

j
L1 (x) = f (x) − λj g (x) − μi (−xi )
j=1 i=1

The necessary conditions for x∗ to solve the problem are

m
∂f (x∗ ) ∂gj (x∗ )
− λj + μi = 0, ∀i = 1, . . . , n
∂xi ∂xi
j=1

λi ≥ 0 and λj = 0 if gj (x∗ ) < 0, ∀j = 1, . . . , m

μi ≥ 0 and μi = 0 if xi > 0, ∀i = 1, . . . , n

94
 To reduce this collection of m + n constraints and m + n Lagrangian multipliers,
the necessary conditions for the optimization problem are sometime formulated slightly
differently below.
m
∂f (x∗ ) ∂gj (x∗ )
− λj ≤ 0 (= 0 if x∗i > 0), ∀i = 1, . . . , n
∂xi ∂xi
j=1

This formulation follows from the first order condition.

m
∂f (x∗ ) ∂gj (x∗ )
− λj = −μi , ∀i = 1, . . . , n
∂xi ∂xi
j=1

Note also that μi ≥ 0 and −μi = 0 if xi > 0

6.7 Concave Programming Problems

The constrained maximization problem is said to be a concave programming program in
the case when f (·) is concave and each gj is a convex function. In this case, the set of
feasible vectors satisfying the m constraints is convex. We write the concave program as
follows:
max f (x) subject to g(x) ≤ 0
where g(x) = (g1 (x), . . . , gm (x)) and 0 = (0, . . . , 0).

Definition 6.3 Let f (·) be concave on convex set S ⊂ Rn , and let x0 be an interior
point in S. Then, there exists a vector p ∈ Rn such that for all x ∈ S,
f (x) − f (x0 ) ≤ p · (x − x0 ).
A vector p that satisfies the above inequality is called a supergradient for f at x0 .

Definition 6.4 The nonlinear programming problem satisfies the Slater qualification
if there exists a vector z ∈ Rn such that g(z)  0, i.e., gj (z) < 0 for all j.

Theorem 6.13 (Necessary Conditions for Concave Programming) Suppose that

the nonlinear programming is a concave programming satisfying the Slater constraint
qualification. Then, the optimal value function f ∗ (c) is defined for (at least) all c ≥ g(z),
and has a super gradient at 0. Furthermore, if λ is any supergradient of f ∗ at 0, then
λ ≥ 0, and any solution x∗ of the concave programming problem is an unconstrained max-
imum point of the Lagrangian L(x, λ) = f (x) − λ · g(x) which also satisfies λ · g(x∗ ) = 0
(the complementary slackness condition).
Proof of Necessity for Concave Programming: We consider only the special
but usual case where, for all c ∈ Rm , the feasible set of points x that satisfy g(x) ≤ c
is bounded, so compact because of the assumption that the functions gj are C 1 , i.e.,
continuous. (See Theorem 3.16). In this case, f ∗ (c) is defined as a maximum value
whenever there exists at least one x satisfying g(x) ≤ c, which is certainly true when
c ≥ g(z). Then, f ∗ is defined for all c ≥ g(z).

95
Theorem 6.14 (Sufficient Conditions for Concave Programming) Consider the
nonlinear programming problem with f (·) concave and g(·) convex, and assume that there
exists a vector λ ≥ 0 and a feasible vector x∗ which together have the property that x∗
maximizes f (x) − λ · g(x) among all x ∈ Rn , and λ · g(x∗ ) = 0. Then, x∗ solves the
original concave problem and λ is a supergradient for f ∗ at 0.

6.8 Quasiconcave Programming

The following theorem is important for economists, because in many economic opti-
mization problems, the objective function is assumed to be quasiconcave, rather than
concave.

Theorem 6.15 (Arrow and Enthoven (1961) in Econometrica ) (Sufficient Con-

ditions for Quasiconcave Programming): Consider the constrained optimization
problem where the objective function f (·) is C 1 and quasiconcave. Assume that there
exist numbers λ1 , . . . , λm and a vector x∗ such that

1. x∗ is feasible and satisfies the Kuhn-Tucker conditions.

2. ∇f (x∗ )
= 0.

3. λj gj (x) is quasiconvex for each j = 1, . . . , m.

Then, x∗ is optimal.

6.9 Appendix: Linear Programming

96
Chapter 7

Differential Equations

7.1 Introduction
What is a differential equation? As the name suggests, it is an equation. Unlike ordinary
algebraic equations, in a differential equation:

• The unknown is a function, not a number.

• The equation includes one or more of the derivatives of the function.

An ordinary differential equation is one for which the unknown is a function of only
one variable. Partial differential equations are equations where the unknown is a function
of two or more variables, and one or more the partial derivatives of the function are
included. In this chapter, I restrict attention to first-order ordinary differential equations
– that is, equations where the first-order derivatives of the unknown functions of one
variable are included.

Consider the following differential equation:

dx
ẋ = = ax (∗)
dt
where x = x(t) is a real-valued function of t ∈ R and ẋ = dx/dt is the first-order
derivative of x(t). The above equation says that for any t ∈ R,

x (t) = ax(t),

where a is some constant. I propose here x(t) = Keat as a solution to the differential
equation.

97
Chapter 8

Fixed Point Theorems

The fixed point problem is

Given a set S ⊂ Rn and a function f : S → S, is there an x ∈ S such that

f (x) = x?

The problem of finding the zeros of a function, f (·), and x ∈ S such that f (x) = 0,
can be converted into a fixed point problem. To see this, observe that f (x) = 0 iff
g(x) = x where g(x) = f (x) + x. The unconstrained optimization with concave objective
function is a special case of the fixed point theorem. In this case, the optimal solution is
found by solving ∇f (x) = 0.

8.1 Banach Fixed Point Theorem

Definition 8.1 A function f : S → S is said to be a contraction mapping if d(f (x), f (y)) ≤
θd(x, y) for all x, y ∈ S, where 0 ≤ θ < 1 is a fixed constant.

Theorem 8.1 (Banach Fixed Point Theorem) Let S ⊂ Rn be closed and f : S → S

a contraction mapping. Then, there exists a unique x ∈ S such that f (x) = x.
Proof of Banach Fixed Point Theorem: We define the norm of vectors in Rn as
follows.

x ≡ max |xi |

1≤n

We use this norm in the proof of the implicit function theorem (Theorem 5.17). Choose
any x0 ∈ S and let xk = f (xk−1 ). If a sequence {xk } has a limit x∗ , then x∗ ∈ S because
S is closed, and f (x∗ ) = x∗ . Therefore, it suffices to prove that {xk } has a limit. We use
the Cauchy criterion. Pick q > p. Then,
6 6
6 q−1 6 q−1
6 6
xq − xp  = 6 6 6
(xk+1 − xk )6 ≤ xk+1 − xk 
6k=p 6 
k=p
Minkowski inequality

98
But

xk+1 − xk  = f (xk ) − f (xk−1 )| ≤ θxk − xk−1 .

Repeated application of the above yields

xk+1 − xk  ≤ θ k x1 − x0 

Hence
q−1

xq − xp  ≤ θ k x1 − x0 
k=p
!
≤ x1 − x0  θ p + θ p+1 + · · ·
θp
= x1 − x0  → 0 as p, q → ∞
1−θ
Because θ p → 0 as p → ∞ due to θ < 1. 

8.2 Brouwer Fixed Point Theorem

Lemma 8.1 If f : [0, 1] → [0, 1] is continuous, there exists x ∈ [0, 1] such that f (x) = x.
Proof of Lemma 8.1: Each x ∈ [0, 1] can be represented as a convex combination
of the end points of the interval:

x = (1 − x) × 0 + x × 1

The same will be true for f (x). So we express each x ∈ [0, 1] as a pair of nonnegative
numbers (x1 , x2 ) = (1 − x, x) that add to one. When expressing f (x) in this way, we will
write it as (f1 (x), f2 (x)) = (1 − f (x), f (x)). Suppose for a contradiction, that f has no
fixed point.

Since f : [0, 1] → [0, 1] we can think of the function f (·) as moving each point x ∈ [0, 1]
either to the right (if f (x) > x) or to the left (if f (x) < x). The assumption that f (·)
has no fixed point eliminates the possibility that f (·) leaves the position of x unchanged.

Given any x ∈ [0, 1], we label it with a “+” if f1 (x) < x1 (move to the right) and
label it “−” if f1 (x) > x1 (move to the left). The assumption of no fixed point implies
f1 (x)
= 1 − x for all x ∈ [0, 1]. Thus, the labeling scheme is well defined. Notice that
the point 0 will be labeled (+) and the point 1 will be labeled (−).

Choose any finite partition, Π0 , of the interval [0, 1] into smaller intervals.

Claim 8.1 The partition Π0 must contain a subinterval [x0 , y 0 ] whose endpoints have
different labels.

99
 Proof of Claim 8.1: Every endpoint of these subintervals is labeled either (+) or
(−). The point “0”, which must be the endpoint of the subinterval of Π0 , has label (+).
The point “1” has label (−). As we travel from 0 to 1 (left to right), we leave a point
labeled (+) and arrive at a point labeled (−). At some point, we must pass through a
subintervals which has endpoints with different labels. 

Now take the partition Π0 and form a new partition Π1 , finer than the first by taking
all the subintervals in Π0 whose endpoints have different labels and cutting them in half.
In Π1 , there must exist at least one subinterval, [x1 , y 1 ] with endpoints having different
labels. Repeat this procedure indefinitely.

This produces an infinite sequence of subintervals {(xk , y k )} shrinking in size with

different labels at the endpoints. Furthermore, we can choose a subsequence of them so
that the left hand endpoint, xk , is labeled (+) and the right hand endpoint y k is labeled
(−). Since these intervals live in [0, 1], their lengths are bounded. Therefore, by Bolzano-
Weierstrass theorem (Theorem 3.13), there is a convergent subsequence of them, with
|xk − y k | → 0 as k → ∞. By continuity of f (·), |f (xk ) − f (y k )| → 0 as k → ∞.

Let z be the limit point of {xk } and {y k }. By continuity, f (xk ) and f (y k ) both
converge to f (z). Since each xk is labeled (+) and each y k is labeled (−), for each k, we
have f1 (xk ) < xk1 and in the limit f1 (z) ≤ z1 . For each k, we have f1 (y k ) > y1k and in
the limit f1 (z) ≥ z1 . Thus, f1 (z) ≤ z1 and f1 (z) ≥ z1 . This implies that f1 (z) = z, i.e.,
f (z) = z, a fixed point. This is a contradiction. 
n
Definition 8.2 The n-simplex is the set Δn = {x ∈ Rn | i=1 xi = 1 and xi ≥ 0 for
all i = 1, . . . , n}.
From the definition, Δn is convex and compact. We also see that this is an (n − 1)-
dimensional object.

Lemma 8.2 If f : Δn → Δn is a continuous function, then there exists x ∈ Δn such

that f (x) = x.
We skip the proof of Lemma 8.2. We should note that Lemma 8.1 is a special case of
Lemma 8.2. Before showing Brouwer fixed point theorem, we need some preliminaries.

Definition 8.3 A set A is topologically equivalent to a set B if there exists a con-

tinuous function g with continuous inverse such that g(A) = B and g−1 (B) = A.
Observe that topological equivalence is a weaker requirement than that for the inverse
function theorem. Do you see why? The closed n-ball of center x0 in Rn is the set
{x ∈ Rn |d(x, x0 ) ≤ 1}. Note that a closed n-ball is of dimension n.

Theorem 8.2 A nonempty compact convex set S ⊂ Rn of dimension m ≤ n is topolog-

ically equivalent to a closed ball in Rm .

100
 We skip the proof of Theorem 8.2. Now, it is time to prove Brouwer fixed point
theorem.

Theorem 8.3 (Brouwer Fixed Point Theorem (1912)) If S ⊂ Rn is compact and

convex and f : S → S is continuous, there exists x ∈ S such that f (x) = x.
Proof of Brouwer Fixed Point Theorem: Lemma 8.2 shows that a continuous
function f : Δn → Δn must have a fixed point. Then, it only remains to prove that there
is no loss of generality to assume that S = Δn as any compact convex set of dimension
n − 1 in Rn . To do so, we make use of the topological equivalence of compact convex
sets.

If S is a compact convex set of dimension n − 1, we know from Theorem 8.2 that

there is a g : S → Δn and g−1 : Δn → S such that g and g−1 are continuous. Define
h : Δn → Δn as follows:
" !#
h(x) = g f g−1 (x) .

Since
" continuous, by Lemma 8.2, it has a fixed point x∗ . Therefore, h(x∗ ) =
h(·) is !#
g f g (x ) = x∗ . We have f (g−1 (x∗ )) = g−1 (x∗ ). So, g−1 (x∗ ) is a fixed point for f .
−1 ∗

101
Chapter 9

Topics on Convex Sets

9.1 Separation Theorems

If a is a nonzero vector in Rn and α is a real number, then the set
H = {x ∈ Rn |a · x = α}
is a hyperplane in Rn , with a as its normal. Moreover, the hyperplane H separates Rn
into two closed half-spaces.

If S and T are subsets of Rn , then H is said to separate S and T if S is contained in

one of the closed half-spaces determined by H and T is contained in the other. In other
words, S and T can be separated by a hyperplane if there exists a vector a
= 0 and a
scalar α such that for all x ∈ S and y ∈ T ,
a·x≤α ≤a·y
If both inequalities are strict, then the hyperplane H = {x ∈ Rn |a · x = α} strictly
separates S and T .

Theorem 9.1 (Strict Separation Theorem) Let S be a closed convex set in Rn , and
let y be a point in Rn that does not belong to S. Then there exists a nonzero vector
a ∈ Rn \{0} and a number α ∈ R such that
a·x<α <a·y
for all x ∈ S. For every such α, the hyperplane H = {x ∈ Rn |a · x = α} strictly separates
S and y.
Proof of Theorem 9.1: Because S is a closed set, among all the points of S, there
is one w = (w1 , . . . , wn ) that is closest to y. For this we can suppose that there is no
such closest point. Then, closedness of S gives us a contradiction. You should fill the
gap in the above argument. Let a = y − w. Since w ∈ S and y ∈ / S, it follows that a
= 0.
Note that a · (y − w) = a · a > 0, and so a · w < a · y. Suppose we prove that
a · x ≤ a · w ∀x ∈ S (∗)

102
Then, the theorem is true for every number α ∈ (a · w, a · y). Now, it remains to show
(∗). Let x be any point in S. Since S is convex, λx + (1 − λ)w ∈ S for each λ ∈ [0, 1].
Now define g(λ) as the square of the distance from λx + (1 − λ)w to the point y:

g(λ) = y − (λx + (1 − λ)w)2 = y − w + λ(w − x)2


Using the chain rule, we obtain g (λ) = 2 (y − w + λ(w − x)) · (w − x). Also g(0) = y −
w2 , the square of the distance between y and w. But w is the point in S that is closest to

y, so g(λ) ≥ g(0) for all λ ∈ [0, 1]. It follows that 0 ≤ g (0) = 2(y−w)·(w−x) = a·(w−x),
which proves (∗). 

In the proof of the above theorem, it was essential that y did not belong to S. If S
is an arbitrary convex set (not necessarily closed), and if y is not an interior point of S,
then it seems plausible that y can still be separated from S by a hyperplane. If y is a
boundary point of S, such a hyperplane is called a supporting hyperplane to S at y.

Theorem 9.2 (Separating Hyperplane) Let S be a convex set in Rn and suppose

y = (y1 , . . . , yn ) is not an interior point of S. Then, there exists a nonzero vector
a ∈ Rn such that

a·x≤a·y for every x ∈ S

Proof of Theorem 9.2: Let S̄ be the closure of S. Because S is convex, so is S̄. Do

you see why? Because y is not an interior point of S and S is convex, y is not an interior
point of S̄. Hence, there is a sequence {yk } of points for which yk ∈ / S̄ for each k and
yk → y as k → ∞. Now, yk ∈ / S̄ and S̄ is closed and convex, so according to the previous
theorem (Theorem 8.4), for each k, there exists a vector ak
= 0 such that ak ·x < ak ·yk for
all x ∈ S̄. Without loss of generality, we can assume that ak  = 1 for each k. Then, {ak }
is a sequence of vectors in {x ∈ Rn |x = 1} which is a compact set. Bolzano-Weierstrass
theorem (Theorem 3.13) shows that {ak } has a convergent subsequence {akj }∞ j=1 . Let
a = limj→∞ akj . Then, a · x = limj→∞ akj · x ≤ limj→∞ akj · ykj = a · y for every x ∈ S̄, as
required. Here we make use of continuity of linear functions. Moreover, we can confirm
that a
= 0 because a = limj→∞ akj  = 1. 

Theorem 9.3 (Separating Hyperplane Theorem) Let S and T be two disjoint nonempty
convex sets in Rn . Then, there exists a nonzero vector a ∈ Rn and a scalar α ∈ R such
that

a·x≤α ≤a·y for all x ∈ S and all y ∈ T

Thus, S and T are separated by the hyperplane H = {z ∈ Rn |a · z = α}.

Proof of Separating Hyperplane Theorem: Let W = S − T be the vector
difference of the two convex sets S and T . Since S and T are disjoint, 0 ∈
/ W . First, I
claim that W is convex.

Claim 9.1 W is convex.

103
  
Proof of Claim 9.1: Let w, w ∈ W . By definition of W , there are s, s ∈ S and
   
t, t ∈ T such that w = s − t and w = s − t . Let α ∈ [0, 1]. What we want to show is

that αw + (1 − α)w ∈ W . We compute the convex combination below.
  
αw + (1 − α)w = αs − αt + (1 − α)s − (1 − α)t
) 
* ) 
*
= αs + (1 − α)s − αt + (1 − α)t
 
Since S and T are convex, αs + (1 − α)s ∈ S and αt + (1 − α)t ∈ T . This implies that

αw + (1 − α)w ∈ S − T = W . 

Hence, by the previous theorem (Theorem 8.5), there exists an a

= 0 such that
a · w ≤ a · 0 = 0 for all w ∈ W . Let x ∈ S and y ∈ T be any two points of these sets.
Then w = x − y ∈ W by definition, so a · (x − y) ≤ 0. Hence
a·x≤a·y for all x ∈ S and all y ∈ T (∗∗)
From (∗∗) it follows that the set A = {a · x|x ∈ S} is bounded above by a · y for any
y ∈ T . By Fact 2.1 (Least Upper Bound Principle), A has a supremum α. Since α
is the least upper bound of A, it follows that α ≤ a · y for every y ∈ T . Therefore,
a · x ≤ α ≤ a · y for all x ∈ S and all y ∈ T . Thus, S and T are separated by the
hyperplane {z ∈ Rn |a · z = α}. 

Theorem 9.4 Let S and T be two disjoint, nonempty, closed, convex sets in Rn with S
being bounded. Then, there exists a nonzero vector a ∈ Rn and a scalar α ∈ R such that
a·x>α >a·y for all x ∈ S and all y ∈ T .

Lemma 9.1 Let A be an m × n matrix, then cone(A) is a convex set.

Lemma 9.2 Let A be an m × n matrix, then cone(A) is a closed set.

Theorem 9.5 (Farkas Lemma) Let A be an m × n matrix, b ∈ Rm and F = {x ∈

Rn | Ax = b, x ≥ 0}. Then either F
= ∅ or there exists y ∈ Rm such that yA ≥ 0 and
yb < 0 but not both.
Proof of Farkas Lemma:

9.2 Polyhedrons and Polytopes

Recall that a set C of vectors is called a cone if λx ∈ C whenever x ∈ C and λ > 0.

Definition 9.1 A cone C ⊂ Rn is polyhedral if there is a matrix A such that C = {x ∈

Rn |Ax ≤ 0}.
Geometrically, a polyhedral cone is the intersection of a finite number of half-spaces
through the origin.

Theorem 9.6 (Farkas-Minkowski-Weyl) A cone C is polyhedral if and only if there

is a finite matrix A such that C = cone(A).

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9.3 Dimension of a Set

9.4 Properties of Convex Sets

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